Split (1)

train · 14.7k rows

id stringlengths 2 8	url stringlengths 31 117	title stringlengths 1 71	text stringlengths 153 118k	topic stringclasses 4 values	section stringlengths 4 49 ⌀	sublist stringclasses 9 values
37515	https://en.wikipedia.org/wiki/Shaped%20charge	Shaped charge	A shaped charge is an explosive charge shaped to focus the effect of the explosive's energy. Different types of shaped charges are used for various purposes such as cutting and forming metal, initiating nuclear weapons, penetrating armor, or perforating wells in the oil and gas industry. A typical modern shaped charge, with a metal liner on the charge cavity, can penetrate armor steel to a depth of seven or more times the diameter of the charge (charge diameters, CD), though depths of 10 CD and above have been achieved. Contrary to a misconception, possibly resulting from the acronym for high-explosive anti-tank, HEAT, the shaped charge does not depend in any way on heating or melting for its effectiveness; that is, the jet from a shaped charge does not melt its way through armor, as its effect is purely kinetic in nature – however the process creates significant heat and often has a significant secondary incendiary effect after penetration. Munroe effect The Munroe or Neumann effect is the focusing of blast energy by a hollow or void cut on a surface of an explosive. The earliest mention of hollow charges were mentioned in 1792. Franz Xaver von Baader (1765–1841) was a German mining engineer at that time; in a mining journal, he advocated a conical space at the forward end of a blasting charge to increase the explosive's effect and thereby save powder. The idea was adopted, for a time, in Norway and in the mines of the Harz mountains of Germany, although the only available explosive at the time was gunpowder, which is not a high explosive and hence incapable of producing the shock wave that the shaped-charge effect requires. The first true hollow charge effect was achieved in 1883, by Max von Foerster (1845–1905), chief of the nitrocellulose factory of Wolff & Co. in Walsrode, Germany. By 1886, Gustav Bloem of Düsseldorf, Germany, had filed for hemispherical cavity metal detonators to concentrate the effect of the explosion in an axial direction. The Munroe effect is named after Charles E. Munroe, who discovered it in 1888. As a civilian chemist working at the U.S. Naval Torpedo Station at Newport, Rhode Island, he noticed that when a block of explosive guncotton with the manufacturer's name stamped into it was detonated next to a metal plate, the lettering was cut into the plate. Conversely, if letters were raised in relief above the surface of the explosive, then the letters on the plate would also be raised above its surface. In 1894, Munroe constructed his first crude shaped charge: Among the experiments made ... was one upon a safe twenty-nine inches cube, with walls four inches and three quarters thick, made up of plates of iron and steel ... When a hollow charge of dynamite nine pounds and a half in weight and untamped was detonated on it, a hole three inches in diameter was blown clear through the wall ... The hollow cartridge was made by tying the sticks of dynamite around a tin can, the open mouth of the latter being placed downward. Although Munroe's experiment with the shaped charge was widely publicized in 1900 in Popular Science Monthly, the importance of the tin can "liner" of the hollow charge remained unrecognized for another 44 years. Part of that 1900 article was reprinted in the February 1945 issue of Popular Science, describing how shaped-charge warheads worked. It was this article that at last revealed to the general public how the United States Army bazooka actually worked against armored vehicles during WWII. In 1910, Egon Neumann of Germany discovered that a block of TNT, which would normally dent a steel plate, punched a hole through it if the explosive had a conical indentation. The military usefulness of Munroe's and Neumann's work was unappreciated for a long time. Between the world wars, academics in several countries Myron Yakovlevich Sukharevskii (Мирон Яковлевич Сухаревский) in the Soviet Union, William H. Payment and Donald Whitley Woodhead in Britain, and Robert Williams Wood in the U.S. recognized that projectiles could form during explosions. In 1932 Franz Rudolf Thomanek, a student of physics at Vienna's Technische Hochschule, conceived an anti-tank round that was based on the hollow charge effect. When the Austrian government showed no interest in pursuing the idea, Thomanek moved to Berlin's Technische Hochschule, where he continued his studies under the ballistics expert Carl Julius Cranz. There in 1935, he and Hellmuth von Huttern developed a prototype anti-tank round. Although the weapon's performance proved disappointing, Thomanek continued his developmental work, collaborating with Hubert Schardin at the Waffeninstitut der Luftwaffe (Air Force Weapons Institute) in Braunschweig. By 1937, Schardin believed that hollow-charge effects were due to the interactions of shock waves. It was during the testing of this idea that, on February 4, 1938, Thomanek conceived the shaped-charge explosive (or Hohlladungs-Auskleidungseffekt (hollow-charge liner effect)). (It was Gustav Adolf Thomer who in 1938 first visualized, by flash radiography, the metallic jet produced by a shaped-charge explosion.) Meanwhile, Henry Hans Mohaupt, a chemical engineer in Switzerland, had independently developed a shaped-charge munition in 1935, which was demonstrated to the Swiss, French, British, and U.S. militaries. During World War II, shaped-charge munitions were developed by Germany (Panzerschreck, Panzerfaust, Panzerwurfmine, Mistel), Britain (No. 68 AT grenade, PIAT, Beehive cratering charge), the Soviet Union (RPG-43, RPG-6), the U.S. (M9 rifle grenade, bazooka), and Italy (Effetto Pronto Speciale shells for various artillery pieces). The development of shaped charges revolutionized anti-tank warfare. Tanks faced a serious vulnerability from a weapon that could be carried by an infantryman or aircraft. One of the earliest uses of shaped charges was by German glider-borne troops against the Belgian Fort Eben-Emael in 1940. These demolition charges – developed by Dr. Wuelfken of the German Ordnance Office – were unlined explosive charges and did not produce a metal jet like the modern HEAT warheads. Due to the lack of metal liner they shook the turrets but they did not destroy them, and other airborne troops were forced to climb on the turrets and smash the gun barrels. Applications Modern military The common term in military terminology for shaped-charge warheads is high-explosive anti-tank (HEAT) warhead. HEAT warheads are frequently used in anti-tank guided missiles, unguided rockets, gun-fired projectiles (both spun (spin stabilized) and unspun), rifle grenades, land mines, bomblets, torpedoes, and various other weapons. Protection During World War II, the precision of the charge's construction and its detonation mode were both inferior to modern warheads. This lower precision caused the jet to curve and to break up at an earlier time and hence at a shorter distance. The resulting dispersion decreased the penetration depth for a given cone diameter and also shortened the optimum standoff distance. Since the charges were less effective at larger standoffs, side and turret skirts (known as Schürzen) fitted to some German tanks to protect against ordinary anti-tank rifles were fortuitously found to give the jet room to disperse and hence also reduce HEAT penetration. The use of add-on spaced armor skirts on armored vehicles may have the opposite effect and actually increase the penetration of some shaped-charge warheads. Due to constraints in the length of the projectile/missile, the built-in stand-off on many warheads is less than the optimum distance. In such cases, the skirting effectively increases the distance between the armor and the target, and the warhead detonates closer to its optimum standoff. Skirting should not be confused with cage armor which is primarily used to damage the fusing system of RPG-7 projectiles, but can also cause a HEAT projectile to pitch up or down on impact, lengthening the penetration path for the shaped charge's penetration stream. If the nose probe strikes one of the cage armor slats, the warhead will function as normal. Non-military In non-military applications shaped charges are used in explosive demolition of buildings and structures, in particular for cutting through metal piles, columns and beams and for boring holes. In steelmaking, small shaped charges are often used to pierce taps that have become plugged with slag. They are also used in quarrying, breaking up ice, breaking log jams, felling trees, and drilling post holes. Shaped charges are used most extensively in the petroleum and natural gas industries, in particular in the completion of oil and gas wells, in which they are detonated to perforate the metal casing of the well at intervals to admit the influx of oil and gas. Another use in the industry is to put out oil and gas fires by depriving the fire of oxygen. A shaped charge was used on the Hayabusa2 mission on asteroid 162173 Ryugu. The spacecraft dropped the explosive device onto the asteroid and detonated it with the spacecraft behind cover. The detonation dug a crater about 10 meters wide, to provide access to a pristine sample of the asteroid. Function A typical device consists of a solid cylinder of explosive with a metal-lined conical hollow in one end and a central detonator, array of detonators, or detonation wave guide at the other end. Explosive energy is released directly away from (normal to) the surface of an explosive, so shaping the explosive will concentrate the explosive energy in the void. If the hollow is properly shaped, usually conically, the enormous pressure generated by the detonation of the explosive drives the liner in the hollow cavity inward to collapse upon its central axis. The resulting collision forms and projects a high-velocity jet of metal particles forward along the axis. Most of the jet material originates from the innermost part of the liner, a layer of about 10% to 20% of the thickness. The rest of the liner forms a slower-moving slug of material, which, because of its appearance, is sometimes called a "carrot". Because of the variation along the liner in its collapse velocity, the jet's velocity also varies along its length, decreasing from the front. This variation in jet velocity stretches it and eventually leads to its break-up into particles. Over time, the particles tend to fall out of alignment, which reduces the depth of penetration at long standoffs. At the apex of the cone, which forms the very front of the jet, the liner does not have time to be fully accelerated before it forms its part of the jet. This results in its small part of jet being projected at a lower velocity than jet formed later behind it. As a result, the initial parts of the jet coalesce to form a pronounced wider tip portion. Most of the jet travels at hypersonic speed. The tip moves at 7 to 14 km/s, the jet tail at a lower velocity (1 to 3 km/s), and the slug at a still lower velocity (less than 1 km/s). The exact velocities depend on the charge's configuration and confinement, explosive type, materials used, and the explosive-initiation mode. At typical velocities, the penetration process generates such enormous pressures that it may be considered hydrodynamic; to a good approximation, the jet and armor may be treated as inviscid, compressible fluids (see, for example,), with their material strengths ignored. A recent technique using magnetic diffusion analysis showed that the temperature of the outer 50% by volume of a copper jet tip while in flight was between 1100K and 1200K, much closer to the melting point of copper (1358 K) than previously assumed. This temperature is consistent with a hydrodynamic calculation that simulated the entire experiment. In comparison, two-color radiometry measurements from the late 1970s indicate lower temperatures for various shaped-charge liner material, cone construction and type of explosive filler. A Comp-B loaded shaped charge with a copper liner and pointed cone apex had a jet tip temperature ranging from 668 K to 863 K over a five shot sampling. Octol-loaded charges with a rounded cone apex generally had higher surface temperatures with an average of 810 K, and the temperature of a tin-lead liner with Comp-B fill averaged 842 K. While the tin-lead jet was determined to be liquid, the copper jets are well below the melting point of copper. However, these temperatures are not completely consistent with evidence that soft recovered copper jet particles show signs of melting at the core while the outer portion remains solid and cannot be equated with bulk temperature. The location of the charge relative to its target is critical for optimum penetration for two reasons. If the charge is detonated too close there is not enough time for the jet to fully develop. But the jet disintegrates and disperses after a relatively short distance, usually well under two meters. At such standoffs, it breaks into particles which tend to tumble and drift off the axis of penetration, so that the successive particles tend to widen rather than deepen the hole. At very long standoffs, velocity is lost to air drag, further degrading penetration. The key to the effectiveness of the hollow charge is its diameter. As the penetration continues through the target, the width of the hole decreases leading to a characteristic "fist to finger" action, where the size of the eventual "finger" is based on the size of the original "fist". In general, shaped charges can penetrate a steel plate as thick as 150% to 700% of their diameter, depending on the charge quality. The figure is for basic steel plate, not for the composite armor, reactive armor, or other types of modern armor. Liner The most common shape of the liner is conical, with an internal apex angle of 40 to 90 degrees. Different apex angles yield different distributions of jet mass and velocity. Small apex angles can result in jet bifurcation, or even in the failure of the jet to form at all; this is attributed to the collapse velocity being above a certain threshold, normally slightly higher than the liner material's bulk sound speed. Other widely used shapes include hemispheres, tulips, trumpets, ellipses, and bi-conics; the various shapes yield jets with different velocity and mass distributions. Liners have been made from many materials, including various metals and glass. The deepest penetrations are achieved with a dense, ductile metal, and a very common choice has been copper. For some modern anti-armor weapons, molybdenum and pseudo-alloys of tungsten filler and copper binder (9:1, thus density is ≈18 Mg/m3) have been adopted. Nearly every common metallic element has been tried, including aluminum, tungsten, tantalum, depleted uranium, lead, tin, cadmium, cobalt, magnesium, titanium, zinc, zirconium, molybdenum, beryllium, nickel, silver, and even gold and platinum. The selection of the material depends on the target to be penetrated; for example, aluminum has been found advantageous for concrete targets. In early antitank weapons, copper was used as a liner material. Later, in the 1970s, it was found tantalum is superior to copper, due to its much higher density and very high ductility at high strain rates. Other high-density metals and alloys tend to have drawbacks in terms of price, toxicity, radioactivity, or lack of ductility. For the deepest penetrations, pure metals yield the best results, because they display the greatest ductility, which delays the breakup of the jet into particles as it stretches. In charges for oil well completion, however, it is essential that a solid slug or "carrot" not be formed, since it would plug the hole just penetrated and interfere with the influx of oil. In the petroleum industry, therefore, liners are generally fabricated by powder metallurgy, often of pseudo-alloys which, if unsintered, yield jets that are composed mainly of dispersed fine metal particles. Unsintered cold pressed liners, however, are not waterproof and tend to be brittle, which makes them easy to damage during handling. Bimetallic liners, usually zinc-lined copper, can be used; during jet formation the zinc layer vaporizes and a slug is not formed; the disadvantage is an increased cost and dependency of jet formation on the quality of bonding the two layers. Low-melting-point (below 500 °C) solder- or braze-like alloys (e.g., Sn50Pb50, Zn97.6Pb1.6, or pure metals like lead, zinc, or cadmium) can be used; these melt before reaching the well casing, and the molten metal does not obstruct the hole. Other alloys, binary eutectics (e.g. Pb88.8Sb11.1, Sn61.9Pd38.1, or Ag71.9Cu28.1), form a metal-matrix composite material with ductile matrix with brittle dendrites; such materials reduce slug formation but are difficult to shape. A metal-matrix composite with discrete inclusions of low-melting material is another option; the inclusions either melt before the jet reaches the well casing, weakening the material, or serve as crack nucleation sites, and the slug breaks up on impact. The dispersion of the second phase can be achieved also with castable alloys (e.g., copper) with a low-melting-point metal insoluble in copper, such as bismuth, 1–5% lithium, or up to 50% (usually 15–30%) lead; the size of inclusions can be adjusted by thermal treatment. Non-homogeneous distribution of the inclusions can also be achieved. Other additives can modify the alloy properties; tin (4–8%), nickel (up to 30% and often together with tin), up to 8% aluminium, phosphorus (forming brittle phosphides) or 1–5% silicon form brittle inclusions serving as crack initiation sites. Up to 30% zinc can be added to lower the material cost and to form additional brittle phases. Oxide glass liners produce jets of low density, therefore yielding less penetration depth. Double-layer liners, with one layer of a less dense but pyrophoric metal (e.g. aluminum or magnesium), can be used to enhance incendiary effects following the armor-piercing action; explosive welding can be used for making those, as then the metal-metal interface is homogeneous, does not contain significant amount of intermetallics, and does not have adverse effects to the formation of the jet. The penetration depth is proportional to the maximum length of the jet, which is a product of the jet tip velocity and time to particulation. The jet tip velocity depends on bulk sound velocity in the liner material, the time to particulation is dependent on the ductility of the material. The maximum achievable jet velocity is roughly 2.34 times the sound velocity in the material. The speed can reach 10 km/s, peaking some 40 microseconds after detonation; the cone tip is subjected to acceleration of about 25 million g. The jet tail reaches about 2–5 km/s. The pressure between the jet tip and the target can reach one terapascal. The immense pressure makes the metal flow like a liquid, though x-ray diffraction has shown the metal stays solid; one of the theories explaining this behavior proposes molten core and solid sheath of the jet. The best materials are face-centered cubic metals, as they are the most ductile, but even graphite and zero-ductility ceramic cones show significant penetration. Explosive charge For optimal penetration, a high explosive with a high detonation velocity and pressure is normally chosen. The most common explosive used in high performance anti-armor warheads is HMX (octogen), although never in its pure form, as it would be too sensitive. It is normally compounded with a few percent of some type of plastic binder, such as in the polymer-bonded explosive (PBX) LX-14, or with another less-sensitive explosive, such as TNT, with which it forms Octol. Other common high-performance explosives are RDX-based compositions, again either as PBXs or mixtures with TNT (to form Composition B and the Cyclotols) or wax (Cyclonites). Some explosives incorporate powdered aluminum to increase their blast and detonation temperature, but this addition generally results in decreased performance of the shaped charge. There has been research into using the very high-performance but sensitive explosive CL-20 in shaped-charge warheads, but, at present, due to its sensitivity, this has been in the form of the PBX composite LX-19 (CL-20 and Estane binder). Other features A 'waveshaper' is a body (typically a disc or cylindrical block) of an inert material (typically solid or foamed plastic, but sometimes metal, perhaps hollow) inserted within the explosive for the purpose of changing the path of the detonation wave. The effect is to modify the collapse of the cone and resulting jet formation, with the intent of increasing penetration performance. Waveshapers are often used to save space; a shorter charge with a waveshaper can achieve the same performance as a longer charge without a waveshaper. Given that the space of possible waveshapes is infinite, machine learning methods have been developed to engineer more optimal waveshapers that can enhance the performance of a shaped charge via computational design. Another useful design feature is sub-calibration, the use of a liner having a smaller diameter (caliber) than the explosive charge. In an ordinary charge, the explosive near the base of the cone is so thin that it is unable to accelerate the adjacent liner to sufficient velocity to form an effective jet. In a sub-calibrated charge, this part of the device is effectively cut off, resulting in a shorter charge with the same performance. Variants There are several forms of shaped charge. Linear shaped charges A linear shaped charge (LSC) has a lining with V-shaped profile and varying length. The lining is surrounded with explosive, the explosive then encased within a suitable material that serves to protect the explosive and to confine (tamp) it on detonation. "At detonation, the focusing of the explosive high pressure wave as it becomes incident to the side wall causes the metal liner of the LSC to collapse–creating the cutting force." The detonation projects into the lining, to form a continuous, knife-like (planar) jet. The jet cuts any material in its path, to a depth depending on the size and materials used in the charge. Generally, the jet penetrates around 1 to 1.2 times the charge width. For the cutting of complex geometries, there are also flexible versions of the linear shaped charge, these with a lead or high-density foam sheathing and a ductile/flexible lining material, which also is often lead. LSCs are commonly used in the cutting of rolled steel joists (RSJ) and other structural targets, such as in the controlled demolition of buildings. LSCs are also used to separate the stages of multistage rockets, and destroy them when they go errant. Explosively formed penetrator The explosively formed penetrator (EFP) is also known as the self-forging fragment (SFF), explosively formed projectile (EFP), self-forging projectile (SEFOP), plate charge, and Misnay-Schardin (MS) charge. An EFP uses the action of the explosive's detonation wave (and to a lesser extent the propulsive effect of its detonation products) to project and deform a plate or dish of ductile metal (such as copper, iron, or tantalum) into a compact high-velocity projectile, commonly called the slug. This slug is projected toward the target at about two kilometers per second. The chief advantage of the EFP over a conventional (e.g., conical) shaped charge is its effectiveness at very great standoffs, equal to hundreds of times the charge's diameter (perhaps a hundred meters for a practical device). The EFP is relatively unaffected by first-generation reactive armor and can travel up to perhaps 1000 charge diameters (CD)s before its velocity becomes ineffective at penetrating armor due to aerodynamic drag, or successfully hitting the target becomes a problem. The impact of a ball or slug EFP normally causes a large-diameter but relatively shallow hole, of, at most, a couple of CDs. If the EFP perforates the armor, spalling and extensive behind armor effects (BAE, also called behind armor damage, BAD) will occur. The BAE is mainly caused by the high-temperature and high-velocity armor and slug fragments being injected into the interior space and the blast overpressure caused by this debris. More modern EFP warhead versions, through the use of advanced initiation modes, can also produce long-rods (stretched slugs), multi-slugs and finned rod/slug projectiles. The long-rods are able to penetrate a much greater depth of armor, at some loss to BAE, multi-slugs are better at defeating light or area targets and the finned projectiles are much more accurate. The use of this warhead type is mainly restricted to lightly armored areas of main battle tanks (MBT) such as the top, belly and rear armored areas. It is well suited for the attack of other less heavily protected armored fighting vehicles (AFV) and in the breaching of material targets (buildings, bunkers, bridge supports, etc.). The newer rod projectiles may be effective against the more heavily armored areas of MBTs. Weapons using the EFP principle have already been used in combat; the "smart" submunitions in the CBU-97 cluster bomb used by the US Air Force and Navy in the 2003 Iraq war employed this principle, and the US Army is reportedly experimenting with precision-guided artillery shells under Project SADARM (Seek And Destroy ARMor). There are also various other projectile (BONUS, DM 642) and rocket submunitions (Motiv-3M, DM 642) and mines (MIFF, TMRP-6) that use EFP principle. Examples of EFP warheads are US patents 5038683 and US6606951. Tandem warhead Some modern anti-tank rockets (RPG-27, RPG-29) and missiles (TOW-2, TOW-2A, Eryx, HOT, MILAN) use a tandem warhead shaped charge, consisting of two separate shaped charges, one in front of the other, typically with some distance between them. TOW-2A was the first to use tandem warheads in the mid-1980s, an aspect of the weapon which the US Army had to reveal under news media and Congressional pressure resulting from the concern that NATO antitank missiles were ineffective against Soviet tanks that were fitted with the new ERA boxes. The Army revealed that a 40 mm precursor shaped-charge warhead was fitted on the tip of the TOW-2 and TOW-2A collapsible probe. Usually, the front charge is somewhat smaller than the rear one, as it is intended primarily to disrupt ERA boxes or tiles. Examples of tandem warheads are US patents 7363862 and US 5561261. The US Hellfire antiarmor missile is one of the few that have accomplished the complex engineering feat of having two shaped charges of the same diameter stacked in one warhead. Recently, a Russian arms firm revealed a 125mm tank cannon round with two same diameter shaped charges one behind the other, but with the back one offset so its penetration stream will not interfere with the front shaped charge's penetration stream. The reasoning behind both the Hellfire and the Russian 125 mm munitions having tandem same diameter warheads is not to increase penetration, but to increase the beyond-armour effect. Voitenko compressor In 1964 a Soviet scientist proposed that a shaped charge originally developed for piercing thick steel armor be adapted to the task of accelerating shock waves. The resulting device, looking a little like a wind tunnel, is called a Voitenko compressor. The Voitenko compressor initially separates a test gas from a shaped charge with a malleable steel plate. When the shaped charge detonates, most of its energy is focused on the steel plate, driving it forward and pushing the test gas ahead of it. Ames Laboratory translated this idea into a self-destroying shock tube. A 66-pound shaped charge accelerated the gas in a 3-cm glass-walled tube 2 meters in length. The velocity of the resulting shock wave was 220,000 feet per second (67 km/s). The apparatus exposed to the detonation was completely destroyed, but not before useful data was extracted. In a typical Voitenko compressor, a shaped charge accelerates hydrogen gas which in turn accelerates a thin disk up to about 40 km/s. A slight modification to the Voitenko compressor concept is a super-compressed detonation, a device that uses a compressible liquid or solid fuel in the steel compression chamber instead of a traditional gas mixture. A further extension of this technology is the explosive diamond anvil cell, utilizing multiple opposed shaped-charge jets projected at a single steel encapsulated fuel, such as hydrogen. The fuels used in these devices, along with the secondary combustion reactions and long blast impulse, produce similar conditions to those encountered in fuel-air and thermobaric explosives. Nuclear shaped charges The proposed Project Orion nuclear propulsion system would have required the development of nuclear shaped charges for reaction acceleration of spacecraft. Shaped-charge effects driven by nuclear explosions have been discussed speculatively, but are not known to have been produced in fact. For example, the early nuclear weapons designer Ted Taylor was quoted as saying, in the context of shaped charges, "A one-kiloton fission device, shaped properly, could make a hole in diameter a thousand feet (305 m) into solid rock." Also, a nuclear driven explosively formed penetrator was apparently proposed for terminal ballistic missile defense in the 1960s.	Technology	Explosive weapons	null
37517	https://en.wikipedia.org/wiki/Anti-tank%20guided%20missile	Anti-tank guided missile	An anti-tank guided missile (ATGM), anti-tank missile, anti-tank guided weapon (ATGW) or anti-armor guided weapon is a guided missile primarily designed to hit and destroy heavily armored military vehicles. ATGMs range in size from shoulder-launched weapons, which can be transported by a single soldier, to larger tripod-mounted weapons, which require a squad or team to transport and fire, to vehicle and aircraft mounted missile systems. Earlier man-portable anti-tank weapons, like anti-tank rifles and magnetic anti-tank mines, generally had very short range, sometimes on the order of metres or tens of metres. Rocket-propelled high-explosive anti-tank (HEAT) systems appeared in World War II and extended range to the order of hundreds of metres, but accuracy was low and hitting targets at these ranges was largely a matter of luck. It was the combination of rocket propulsion and remote wire guidance that made the ATGM much more effective than these earlier weapons, and gave light infantry real capability on the battlefield against post-war tank designs. The introduction of semi-automatic guidance in the 1960s further improved the performance of ATGMs. ATGMs were used by over 130 countries and many non-state actors around the world. Post-Cold-War main battle tanks (MBTs) using composite and reactive armors have proven to be resistant to smaller ATGMs. History World War II Germany developed a design for a wire-guided anti tank missile derived from the Ruhrstahl X-4 air to air missile concept in the closing years of World War II. Known as the X-7, it was probably never used in combat and allegedly had serious guidance to target issues. It never entered service, though a few were produced. Early Cold War: first generation ATGMs First-generation ATGMs use a type of command guidance termed manual command to line of sight (MCLOS). This requires continuous input from an operator using a joystick or similar control system to steer the missile to a target. One disadvantage of this is that an operator must keep the sight's reticle cross hairs on a target and then steer the missile into the cross hairs, i.e., the line-of-sight. To do this, an operator must be well trained (spending many hours on a simulator) and must remain stationary and in view of a target during the flight time of the missile. Because of this, the operator is vulnerable while guiding the missile. In addition to the low kill probability, other problems with first generation ATGMs include slow missile speed, high minimum effective range, and an inability to use top attack missiles. The first system to become operational and to see combat was the French Nord SS.10 during the early 1950s. It entered service in the French Army in 1955. It was also the first anti-tank missile used by the US Army and Israeli Defense Forces. The Malkara missile (named from an Australian Aborigine word for "shield") was another of the earliest ATGMs. It was jointly developed by Australia and the United Kingdom between 1951 and 1954, and was in service from 1958 until gradually replaced by the Vickers Vigilant missile in the late 1960s. It was intended to be light enough to deploy with airborne forces, yet powerful enough to knock out any tank then in service. It used a high-explosive squash head (HESH) warhead. Other early first generation ATGMs include the West German Cobra and the Soviet 9M14 Malyutka. In 2012, first-generation systems were described as obsolete due to low hit probability, a limited ability to penetrate modern armour, and other issues. Still, many countries maintain significant stockpiles. Approximately, first generation ATGMs have an effective range of 1500m and the ability to penetrate 500mm of rolled homogeneous armor. Late Cold War: second generation ATGMs Second-generation semi-automatically command guided to line-of-sight, or semi-automatic command to line of sight (SACLOS) missiles require an operator to only keep the sights on the target until impact. Automatic guidance commands are sent to the missile through wires or radio, or the missile relies on laser marking or a TV camera view from the nose of the missile. Examples are the Russian 9M133 Kornet, Israeli LAHAT, the NLOS version of Spike, and the American Hellfire I missiles. The operator must remain stationary during the missile's flight. The most widely used ATGM of all time, the American BGM-71 TOW, with hundreds of thousands of missiles built, is a second-generation system. Second generation ATGMs are significantly easier to use than first generation systems, and accuracy rates may exceed 90%. Generally they have an effective range of between 2,500 and 5,500 meters and penetration of up to 900 mm of armor. Cost is around $10,000 USD per missile. Post Cold War: third generation ATGMs and later Third-generation "fire-and-forget" missiles rely on a laser, electro-optical imager (IIR) seeker or a W band radar seeker in the nose of the missile. Once the target is identified, the missile needs no further guidance during flight; it is "fire-and-forget", and the missile operator is free to retreat. However, fire-and-forget missiles are more subject to electronic countermeasures than MCLOS and SACLOS missiles. Examples include the German PARS 3 LR and the Israeli Spike. Most modern ATGMs have shaped charge HEAT warheads, designed specifically for penetrating tank armor. Tandem-charge missiles attempt to defeat explosive reactive armour (ERA): the small initial charge sets off the ERA while the follow-up main charge attempts to penetrate the main armor. Top-attack weapons such as the US Javelin, the Swedish Bill and the Indian Nag and MPATGM are designed to strike vehicles from above, where their armor is usually much weaker. Third generation systems and beyond are generally much more expensive than second generation systems. Fourth generation ATGMs Fourth generation fire-and-forget anti tank guided missiles have larger range and rely on a combination of seeker for guidance. Examples include India's SANT, which has a stand-off range of , uses dual seeker configuration of electro-optical thermal imager (EO/IR) and millimeter-wave active radar homing for control and guidance with lock-on before launch and lock-on after launch capabilities. Fifth generation ATGMs Some ATGMs, notably the French Akeron MP and the latest variants of the Israeli Spike (such as the Spike LR2 and ER2), have been called "5th generation" by their manufacturers and marketed as such. They appear to have the following additional or amplified attributes: passive dual-band seeker (TV and uncooled IR); multipurpose tandem warhead; smokeless propellant; less collateral damage; possible counter-active protection system (CAPS) capability; man in the loop technology; emphasis on targets other than tanks; other updates such as artificial intelligence for the missile. Countermeasures Countermeasures against ATGMs include newer armors such as spaced, perforated, composite or explosive reactive armor, jammers like the Russian Shtora, active protection systems (APS) like the Israeli Trophy and the Russian Arena, and other methods. Newer armor Armor systems have continued in development alongside ATGMs, and the most recent generations of armor are specifically tested to be effective against ATGM strikes, either by deforming the missile warhead or fusing to prevent proper detonation (such as in slat armor) or using some form of reactive armor to 'attack' the missile upon impact, disrupting the shaped charge that makes the warhead effective. Both come with the downside of significant weight and bulk. Reactive armor works best when a vehicle is specifically designed with the system integrated and while developments continue to make armor lighter, any vehicle that includes such a system necessitates a powerful engine and often will still be relatively slow. Inclusion of such armor in older vehicles as a part of a re-design is possible, as in the numerous types derived from the T-72. Slat armor is lighter and as such can be added to many vehicles after construction but still adds both bulk and weight. Particularly for vehicles that are designed to be transported by cargo aircraft, slat armor has to be fitted in the field after deployment. Either approach can never offer complete coverage over the vehicle, leaving tracks or wheels particularly vulnerable to attack. Jamming Jamming is potentially an effective countermeasure to specific missiles that are radar guided, however, as a general purpose defense, it is of no use against unguided anti-tank weapons, and as such it is almost never the only defense. If jamming is used continually, it can be extremely difficult for a missile to acquire the target, locking on to the much larger return from the jammer, with the operator unlikely noticing the difference without a radar screen to see the return. However, any missile that has a backup tracking system can defeat jamming. Active Active protection systems show a great deal of promise, both in counteracting ATGMs and unguided weapons. Compared to armor systems, they are very lightweight, can be fitted to almost any vehicle with the internal space for the control system and could, in the future, be a near-perfect defense against any missiles. The weaknesses of the systems include potential developments in missile design such as radar or IR decoys, which would drastically reduce their chance to intercept a missile, as well as technical challenges such as dealing with multiple missiles at once and designing a system that can cover a vehicle from any angle of attack. While these may be answered and allow for lightweight, highly maneuverable vehicles that are strongly defended against missiles and rockets that are extremely well suited for urban and guerrilla warfare, such a system is unlikely to be as effective against kinetic energy projectiles, making it a poor choice for fighting against tanks. As kinetic energy projectiles move faster than guided missiles, this often means that the sensors attached to an active protection system can not keep up. Other Traditionally, before "fire-and-forget" ATGMs were used, the most effective countermeasure was to open fire at the location where the missile was fired from, to either kill the operator or force them to take cover, thus sending the missile off course. Smoke screens can also be deployed from an MBT's smoke discharger, and used to obscure an ATGM operator's line of sight. Other improvised methods used by the Israel Defense Forces to defeat the Saggers involved firing in front of the tank to create dust. While fire-and-forget missiles have definitive advantages in terms of guidance and operator safety, and include abilities such as top attack mode, older missiles continue in use, both in the front line armies of less developed countries, and in reserve service the world over, due to their lower cost or existing stockpiles of less advanced weapons.	Technology	Missiles	null
37526	https://en.wikipedia.org/wiki/433%20Eros	433 Eros	433 Eros is a stony asteroid of the Amor group, and the first discovered, and second-largest near-Earth object. It has an elongated shape and a volume-equivalent diameter of approximately . Visited by the NEAR Shoemaker space probe in 1998, it became the first asteroid ever studied from its own orbit. The asteroid was discovered by German astronomer C. G. Witt at the Berlin Observatory on 13 August 1898 in an eccentric orbit between Mars and Earth. It was later named after Eros, a god from Greek mythology, the son of Aphrodite. He is identified with the planet Venus. History Discovery Eros was discovered on 13 August 1898 by Carl Gustav Witt at Berlin Urania Observatory and Auguste Charlois at Nice Observatory and temporarily labeled D.Q. Witt was taking a two-hour exposure of beta Aquarii to secure astrometric positions of asteroid 185 Eunike. Name Eros is named after the Greek god of love, Erōs. It was the first minor planet to be given a male name; the break with earlier tradition was made because it was the first near-Earth asteroid discovered. Later studies During the opposition of 1900–1901, a worldwide program was launched to make parallax measurements of Eros to determine the solar parallax (or distance to the Sun), with the results published in 1910 by Arthur Hinks of Cambridge and Charles D. Perrine of the Lick Observatory, University of California. Perrine published progress reports in 1906 and 1908. He took 965 photographs with the Crossley Reflector and selected 525 for measurement. A similar program was then carried out, during a closer approach, in 1930–1931 by Harold Spencer Jones. The value of the Astronomical Unit (roughly the Earth-Sun distance) obtained by this program was considered definitive until 1968, when radar and dynamical parallax methods started producing more precise measurements. Eros was the first asteroid detected by the Arecibo Observatory's radar system. Eros was one of the first asteroids visited by a spacecraft, the first one orbited, and the first one soft-landed on. NASA spacecraft NEAR Shoemaker entered orbit around Eros in 2000, and landed in 2001. Mars-crosser Eros is a Mars-crosser asteroid, the first known to come within the orbit of Mars. Objects in such an orbit can remain there for only a few hundred million years before the orbit is perturbed by gravitational interactions. Dynamical system modeling suggests that Eros may evolve into an Earth-crosser within as short an interval as two million years, and has a roughly 50% chance of doing so over a time scale of ~ years. It is a potential Earth impactor, about five times larger than the impactor that created Chicxulub crater and led to the extinction of the non-avian dinosaurs. NEAR Shoemaker survey and landing The NEAR Shoemaker probe visited Eros twice, first with a brief flyby in 1998, and then by orbiting it in 2000, when it extensively photographed its surface. On 12 February 2001, at the end of its mission, it landed on the asteroid's surface using its maneuvering jets. This was the first time a Near Earth asteroid was closely visited by a spacecraft. Physical characteristics Surface gravity depends on the distance from a spot on the surface to the center of a body's mass. Eros's surface gravity varies greatly because Eros is not a sphere but an elongated peanut-shaped object. The daytime temperature on Eros can reach about at perihelion. Nighttime measurements fall near . Eros's density is 2.67 g/cm3, about the same as the density of Earth's crust. NEAR scientists have found that most of the larger rocks strewn across Eros were ejected from a single crater in an impact approximately 1 billion years ago. (The crater involved was proposed to be named "Shoemaker", but is not recognized as such by the International Astronomical Union (IAU), and has been formally designated Charlois Regio.) This event may also be responsible for the 40 percent of the Erotian surface that is devoid of craters smaller than 0.5 kilometers across. It was originally thought that the debris thrown up by the collision filled in the smaller craters. An analysis of crater densities over the surface indicates that the areas with lower crater density are within 9 kilometers of the impact point. Some of the lower density areas were found on the opposite side of the asteroid but still within 9 kilometers. It is thought that seismic shockwaves propagate through the asteroid, shaking smaller craters into rubble. Since Eros is irregularly shaped, parts of the surface antipodal to the point of impact can be within 9 kilometres of the impact point (measured in a straight line through the asteroid) even though some intervening parts of the surface are more than 9 kilometres away in straight-line distance. A suitable analogy would be the distance from the top centre of a bun to the bottom centre as compared to the distance from the top centre to a point on the bun's circumference: top-to-bottom is a longer distance than top-to-periphery when measured along the surface but shorter than it in direct straight-line terms. Compression from the same impact is believed to have created the thrust fault Hinks Dorsum. A phenomenon named dust ponds were discovered in the asteroid in October 2000. Dust ponds are a phenomenon where pockets of dust are seen in airless celestial bodies. These are smooth deposits of dust accumulated in depressions on the surface of the body (like craters), contrasting from the rocky terrain around them. They typically have different color and albedo compared to the surrounding areas. The asteroid contains lots of large craters more than 200 m in diameter. Their number is near to the saturation point of these craters. But craters smaller than that are relatively low. Suggesting that some process of erasure has covered them up. The floors of some craters are covered with smooth and flat areas (less than 10° slope). Such dust ponds are characterized by slightly bluer colour compared to the surrounding terrain. 334 of such ponds are identified, with a diameter of 10m. 255 of these are larger than 30m, and 231 (or 91%) are found within 30° from equator. Data from the Near Earth Asteroid Rendezvous spacecraft collected on Eros in December 1998 suggests that it could contain 20 billion tonnes of aluminum and similar amounts of metals that are rare on Earth, such as gold and platinum. Visibility from Earth On 31 January 2012, Eros passed Earth at , about 70 times the distance to the Moon, with a visual magnitude of +8.1. During rare oppositions, every 81 years, such as in 1975 and 2056, Eros can reach a magnitude of +7.0, which is brighter than Neptune and brighter than any main-belt asteroid except 1 Ceres, 4 Vesta and, rarely, 2 Pallas and 7 Iris. Under this condition, the asteroid actually appears to stop, but unlike the normal condition for a body in heliocentric conjunction with Earth, its retrograde motion is very small. For example, in January and February 2137, it moves retrograde only 34 minutes in right ascension. In popular culture In the novel and television series The Expanse, a catastrophic science experiment is unleashed on a civilian population living within tunnels cut through Eros. This so-called "Eros Incident" ends with the asteroid mysteriously breaking its usual orbit and crashing into Venus. It makes an appearance in the novel (and its film adaptation) Ender's Game by Orson Scott Card, serving as a base for humanity and the location of Command School after having been captured from the invading aliens (the Formics) prior to the initial novel who had used the asteroid as their forward operating base in their previous invasion. In the Space Angel episode 'Visitors from Outer Space' (title text not quite matching narration), Scott McCloud and his crew are forced to destroy Eros by deflecting it into the Sun, after it becomes a hazard to spacecraft navigation. It is the setting for the entirety of the plot of the novel Captive Universe by Harry Harrison. Gallery	Physical sciences	Solar System	Astronomy
37530	https://en.wikipedia.org/wiki/Tornado	Tornado	A tornado is a violently rotating column of air that is in contact with both the surface of the Earth and a cumulonimbus cloud or, in rare cases, the base of a cumulus cloud. It is often referred to as a twister, whirlwind or cyclone, although the word cyclone is used in meteorology to name a weather system with a low-pressure area in the center around which, from an observer looking down toward the surface of the Earth, winds blow counterclockwise in the Northern Hemisphere and clockwise in the Southern. Tornadoes come in many shapes and sizes, and they are often (but not always) visible in the form of a condensation funnel originating from the base of a cumulonimbus cloud, with a cloud of rotating debris and dust beneath it. Most tornadoes have wind speeds less than , are about across, and travel several kilometers (a few miles) before dissipating. The most extreme tornadoes can attain wind speeds of more than , can be more than in diameter, and can stay on the ground for more than . Various types of tornadoes include the multiple-vortex tornado, landspout, and waterspout. Waterspouts are characterized by a spiraling funnel-shaped wind current, connecting to a large cumulus or cumulonimbus cloud. They are generally classified as non-supercellular tornadoes that develop over bodies of water, but there is disagreement over whether to classify them as true tornadoes. These spiraling columns of air frequently develop in tropical areas close to the equator and are less common at high latitudes. Other tornado-like phenomena that exist in nature include the gustnado, dust devil, fire whirl, and steam devil. Tornadoes occur most frequently in North America (particularly in central and southeastern regions of the United States colloquially known as Tornado Alley; the United States has by far the most tornadoes of any country in the world). Tornadoes also occur in South Africa, much of Europe (except most of the Alps), western and eastern Australia, New Zealand, Bangladesh and adjacent eastern India, Japan, the Philippines, and southeastern South America (Uruguay and Argentina). Tornadoes can be detected before or as they occur through the use of pulse-Doppler radar by recognizing patterns in velocity and reflectivity data, such as hook echoes or debris balls, as well as through the efforts of storm spotters. Tornado rating scales There are several scales for rating the strength of tornadoes. The Fujita scale rates tornadoes by damage caused and has been replaced in some countries by the updated Enhanced Fujita Scale. An F0 or EF0 tornado, the weakest category, damages trees, but not substantial structures. An F5 or EF5 tornado, the strongest category, rips buildings off their foundations and can deform large skyscrapers. The similar TORRO scale ranges from T0 for extremely weak tornadoes to T11 for the most powerful known tornadoes. The International Fujita scale is also used to rate the intensity of tornadoes and other wind events based on the severity of the damage they cause. Doppler radar data, photogrammetry, and ground swirl patterns (trochoidal marks) may also be analyzed to determine intensity and assign a rating. Etymology The word tornado comes from the Spanish (meaning 'thunderstorm', past participle of tronar 'to thunder', itself in turn from the Latin tonāre 'to thunder'). The metathesis of the r and o in the English spelling was influenced by the Spanish tornado (past participle of tornar 'to twist, turn,', from Latin tornō 'to turn'). The English word has been reborrowed into Spanish, referring to the same weather phenomenon. Tornadoes' opposite phenomena are the widespread, straight-line derechos (, from , 'straight'). A tornado is also commonly referred to as a "twister" or the old-fashioned colloquial term cyclone. Definitions A tornado is a violently rotating column of air, in contact with the ground, either pendant from a cumuliform cloud or underneath a cumuliform cloud, and often (but not always) visible as a funnel cloud. For a vortex to be classified as a tornado, it must be in contact with both the ground and the cloud base. The term is not precisely defined; for example, there is disagreement as to whether separate touchdowns of the same funnel constitute separate tornadoes. Tornado refers to the vortex of wind, not the condensation cloud. Funnel cloud A tornado is not necessarily visible; however, the intense low pressure caused by the high wind speeds (as described by Bernoulli's principle) and rapid rotation (due to cyclostrophic balance) usually cause water vapor in the air to condense into cloud droplets due to adiabatic cooling. This results in the formation of a visible funnel cloud or condensation funnel. There is some disagreement over the definition of a funnel cloud and a condensation funnel. According to the Glossary of Meteorology, a funnel cloud is any rotating cloud pendant from a cumulus or cumulonimbus, and thus most tornadoes are included under this definition. Among many meteorologists, the "funnel cloud" term is strictly defined as a rotating cloud which is not associated with strong winds at the surface, and condensation funnel is a broad term for any rotating cloud below a cumuliform cloud. Tornadoes often begin as funnel clouds with no associated strong winds at the surface, and not all funnel clouds evolve into tornadoes. Most tornadoes produce strong winds at the surface while the visible funnel is still above the ground, so it is difficult to discern the difference between a funnel cloud and a tornado from a distance. Outbreaks and families Occasionally, a single storm will produce more than one tornado, either simultaneously or in succession. Multiple tornadoes produced by the same storm cell are referred to as a "tornado family". Several tornadoes are sometimes spawned from the same large-scale storm system. If there is no break in activity, this is considered a tornado outbreak (although the term "tornado outbreak" has various definitions). A period of several successive days with tornado outbreaks in the same general area (spawned by multiple weather systems) is a tornado outbreak sequence, occasionally called an extended tornado outbreak. Characteristics Size and shape Most tornadoes take on the appearance of a narrow funnel, a few hundred meters (yards) across, with a small cloud of debris near the ground. Tornadoes may be obscured completely by rain or dust. These tornadoes are especially dangerous, as even experienced meteorologists might not see them. Small, relatively weak landspouts may be visible only as a small swirl of dust on the ground. Although the condensation funnel may not extend all the way to the ground, if associated surface winds are greater than , the circulation is considered a tornado. A tornado with a nearly cylindrical profile and relatively low height is sometimes referred to as a "stovepipe" tornado. Large tornadoes which appear at least as wide as their cloud-to-ground height can look like large wedges stuck into the ground, and so are known as "wedge tornadoes" or "wedges". The "stovepipe" classification is also used for this type of tornado if it otherwise fits that profile. A wedge can be so wide that it appears to be a block of dark clouds, wider than the distance from the cloud base to the ground. Even experienced storm observers may not be able to tell the difference between a low-hanging cloud and a wedge tornado from a distance. Many, but not all major tornadoes are wedges. Tornadoes in the dissipating stage can resemble narrow tubes or ropes, and often curl or twist into complex shapes. These tornadoes are said to be "roping out", or becoming a "rope tornado". When they rope out, the length of their funnel increases, which forces the winds within the funnel to weaken due to conservation of angular momentum. Multiple-vortex tornadoes can appear as a family of swirls circling a common center, or they may be completely obscured by condensation, dust, and debris, appearing to be a single funnel. In the United States, tornadoes are around across on average. However, there is a wide range of tornado sizes. Weak tornadoes, or strong yet dissipating tornadoes, can be exceedingly narrow, sometimes only a few feet or couple meters across. One tornado was reported to have a damage path only long. On the other end of the spectrum, wedge tornadoes can have a damage path a mile (1.6 km) wide or more. A tornado that affected Hallam, Nebraska on May 22, 2004, was up to wide at the ground, and a tornado in El Reno, Oklahoma on May 31, 2013, was approximately wide, the widest on record. Track length In the United States, the average tornado travels on the ground for . However, tornadoes are capable of both much shorter and much longer damage paths: one tornado was reported to have a damage path only long, while the record-holding tornado for path length—the Tri-State Tornado, which affected parts of Missouri, Illinois, and Indiana on March 18, 1925—was on the ground continuously for . Many tornadoes which appear to have path lengths of or longer are composed of a family of tornadoes which have formed in quick succession; however, there is no substantial evidence that this occurred in the case of the Tri-State Tornado. A 2007 reanalysis of the path suggests that the tornado may have begun further west than previously thought. Appearance Tornadoes can have a wide range of colors, depending on the environment in which they form. Those that form in dry environments can be nearly invisible, marked only by swirling debris at the base of the funnel. Condensation funnels that pick up little or no debris can be gray to white. While traveling over a body of water (as a waterspout), tornadoes can turn white or even blue. Slow-moving funnels, which ingest a considerable amount of debris and dirt, are usually darker, taking on the color of debris. Tornadoes in the Great Plains can turn red because of the reddish tint of the soil, and tornadoes in mountainous areas can travel over snow-covered ground, turning white. Lighting conditions are a major factor in the appearance of a tornado. A tornado which is "back-lit" (viewed with the sun behind it) appears very dark. The same tornado, viewed with the sun at the observer's back, may appear gray or brilliant white. Tornadoes which occur near the time of sunset can be many different colors, appearing in hues of yellow, orange, and pink. Dust kicked up by the winds of the parent thunderstorm, heavy rain and hail, and the darkness of night are all factors that can reduce the visibility of tornadoes. Tornadoes occurring in these conditions are especially dangerous, since only weather radar observations, or possibly the sound of an approaching tornado, serve as any warning to those in the storm's path. Most significant tornadoes form under the storm's updraft base, which is rain-free, making them visible. Also, most tornadoes occur in the late afternoon, when the bright sun can penetrate even the thickest clouds. There is mounting evidence, including Doppler on Wheels mobile radar images and eyewitness accounts, that most tornadoes have a clear, calm center with extremely low pressure, akin to the eye of tropical cyclones. Lightning is said to be the source of illumination for those who claim to have seen the interior of a tornado. Rotation Tornadoes normally rotate cyclonically (when viewed from above, this is counterclockwise in the northern hemisphere and clockwise in the southern). While large-scale storms always rotate cyclonically due to the Coriolis effect, thunderstorms and tornadoes are so small that the direct influence of the Coriolis effect is negligible, as indicated by their large Rossby numbers. Supercells and tornadoes rotate cyclonically in numerical simulations even when the Coriolis effect is neglected. Low-level mesocyclones and tornadoes owe their rotation to complex processes within the supercell and ambient environment. Approximately 1 percent of tornadoes rotate in an anticyclonic direction in the northern hemisphere. Typically, systems as weak as landspouts and gustnadoes can rotate anticyclonically, and usually only those which form on the anticyclonic shear side of the descending rear flank downdraft (RFD) in a cyclonic supercell. On rare occasions, anticyclonic tornadoes form in association with the mesoanticyclone of an anticyclonic supercell, in the same manner as the typical cyclonic tornado, or as a companion tornado either as a satellite tornado or associated with anticyclonic eddies within a supercell. Sound and seismology Tornadoes emit widely on the acoustics spectrum and the sounds are caused by multiple mechanisms. Various sounds of tornadoes have been reported, mostly related to familiar sounds for the witness and generally some variation of a whooshing roar. Popularly reported sounds include a freight train, rushing rapids or waterfall, a nearby jet engine, or combinations of these. Many tornadoes are not audible from much distance; the nature of and the propagation distance of the audible sound depends on atmospheric conditions and topography. The winds of the tornado vortex and of constituent turbulent eddies, as well as airflow interaction with the surface and debris, contribute to the sounds. Funnel clouds also produce sounds. Funnel clouds and small tornadoes are reported as whistling, whining, humming, or the buzzing of innumerable bees or electricity, or more or less harmonic, whereas many tornadoes are reported as a continuous, deep rumbling, or an irregular sound of "noise". Since many tornadoes are audible only when very near, sound is not to be thought of as a reliable warning signal for a tornado. Tornadoes are also not the only source of such sounds in severe thunderstorms; any strong, damaging wind, a severe hail volley, or continuous thunder in a thunderstorm may produce a roaring sound. Tornadoes also produce identifiable inaudible infrasonic signatures. Unlike audible signatures, tornadic signatures have been isolated; due to the long-distance propagation of low-frequency sound, efforts are ongoing to develop tornado prediction and detection devices with additional value in understanding tornado morphology, dynamics, and creation. Tornadoes also produce a detectable seismic signature, and research continues on isolating it and understanding the process. Electromagnetic, lightning, and other effects Tornadoes emit on the electromagnetic spectrum, with sferics and E-field effects detected. There are observed correlations between tornadoes and patterns of lightning. Tornadic storms do not produce more lightning than other storms and some tornadic cells never produce lightning at all. More often than not, overall cloud-to-ground (CG) lightning activity decreases as a tornado touches the surface and returns to the baseline level when the tornado dissipates. In many cases, intense tornadoes and thunderstorms exhibit an increased and anomalous dominance of positive polarity CG discharges. Luminosity has been reported in the past and is probably due to misidentification of external light sources such as lightning, city lights, and power flashes from broken lines, as internal sources are now uncommonly reported and are not known to ever have been recorded. In addition to winds, tornadoes also exhibit changes in atmospheric variables such as temperature, moisture, and atmospheric pressure. For example, on June 24, 2003, near Manchester, South Dakota, a probe measured a pressure decrease. The pressure dropped gradually as the vortex approached then dropped extremely rapidly to in the core of the violent tornado before rising rapidly as the vortex moved away, resulting in a V-shape pressure trace. Temperature tends to decrease and moisture content to increase in the immediate vicinity of a tornado. Life cycle Supercell relationship Tornadoes often develop from a class of thunderstorms known as supercells. Supercells contain mesocyclones, an area of organized rotation a few kilometers/miles up in the atmosphere, usually across. Most intense tornadoes (EF3 to EF5 on the Enhanced Fujita Scale) develop from supercells. In addition to tornadoes, very heavy rain, frequent lightning, strong wind gusts, and hail are common in such storms. Most tornadoes from supercells follow a recognizable life cycle which begins when increasing rainfall drags with it an area of quickly descending air known as the rear flank downdraft (RFD). This downdraft accelerates as it approaches the ground, and drags the supercell's rotating mesocyclone towards the ground with it. Formation As the mesocyclone lowers below the cloud base, it begins to take in cool, moist air from the downdraft region of the storm. The convergence of warm air in the updraft and cool air causes a rotating wall cloud to form. The RFD also focuses the mesocyclone's base, causing it to draw air from a smaller and smaller area on the ground. As the updraft intensifies, it creates an area of low pressure at the surface. This pulls the focused mesocyclone down, in the form of a visible condensation funnel. As the funnel descends, the RFD also reaches the ground, fanning outward and creating a gust front that can cause severe damage a considerable distance from the tornado. Usually, the funnel cloud begins causing damage on the ground (becoming a tornado) within a few minutes of the RFD reaching the ground. Many other aspects of tornado formation (such as why some storms form tornadoes while others do not, or what precise role downdrafts, temperature, and moisture play in tornado formation) are still poorly understood. Maturity Initially, the tornado has a good source of warm, moist air flowing inward to power it, and it grows until it reaches the "mature stage". This can last from a few minutes to more than an hour, and during that time a tornado often causes the most damage, and in rare cases can be more than across. The low pressured atmosphere at the base of the tornado is essential to the endurance of the system. Meanwhile, the RFD, now an area of cool surface winds, begins to wrap around the tornado, cutting off the inflow of warm air which previously fed the tornado. The flow inside the funnel of the tornado is downward, supplying water vapor from the cloud above. This is contrary to the upward flow inside hurricanes, supplying water vapor from the warm ocean below. Therefore, the energy of the tornado is supplied from the cloud above. Dissipation As the RFD completely wraps around and chokes off the tornado's air supply, the vortex begins to weaken, becoming thin and rope-like. This is the "dissipating stage", often lasting no more than a few minutes, after which the tornado ends. During this stage, the shape of the tornado becomes highly influenced by the winds of the parent storm, and can be blown into fantastic patterns. Even though the tornado is dissipating, it is still capable of causing damage. The storm is contracting into a rope-like tube and, due to conservation of angular momentum, winds can increase at this point. As the tornado enters the dissipating stage, its associated mesocyclone often weakens as well, as the rear flank downdraft cuts off the inflow powering it. Sometimes, in intense supercells, tornadoes can develop cyclically. As the first mesocyclone and associated tornado dissipate, the storm's inflow may be concentrated into a new area closer to the center of the storm and possibly feed a new mesocyclone. If a new mesocyclone develops, the cycle may start again, producing one or more new tornadoes. Occasionally, the old (occluded) mesocyclone and the new mesocyclone produce a tornado at the same time. Although this is a widely accepted theory for how most tornadoes form, live, and die, it does not explain the formation of smaller tornadoes, such as landspouts, long-lived tornadoes, or tornadoes with multiple vortices. These each have different mechanisms which influence their development—however, most tornadoes follow a pattern similar to this one. Types Multiple vortex A multiple-vortex tornado is a type of tornado in which two or more columns of spinning air rotate about their own axes and at the same time revolve around a common center. A multi-vortex structure can occur in almost any circulation, but is very often observed in intense tornadoes. These vortices often create small areas of heavier damage along the main tornado path. This is a phenomenon that is distinct from a satellite tornado, which is a smaller tornado that forms very near a large, strong tornado contained within the same mesocyclone. The satellite tornado may appear to "orbit" the larger tornado (hence the name), giving the appearance of one, large multi-vortex tornado. However, a satellite tornado is a distinct circulation, and is much smaller than the main funnel. Waterspout A waterspout is defined by the National Weather Service as a tornado over water. However, researchers typically distinguish "fair weather" waterspouts from tornadic (i.e. associated with a mesocyclone) waterspouts. Fair weather waterspouts are less severe but far more common, and are similar to dust devils and landspouts. They form at the bases of cumulus congestus clouds over tropical and subtropical waters. They have relatively weak winds, smooth laminar walls, and typically travel very slowly. They occur most commonly in the Florida Keys and in the northern Adriatic Sea. In contrast, tornadic waterspouts are stronger tornadoes over water. They form over water similarly to mesocyclonic tornadoes, or are stronger tornadoes which cross over water. Since they form from severe thunderstorms and can be far more intense, faster, and longer-lived than fair weather waterspouts, they are more dangerous. In official tornado statistics, waterspouts are generally not counted unless they affect land, though some European weather agencies count waterspouts and tornadoes together. Landspout A landspout, or dust-tube tornado, is a tornado not associated with a mesocyclone. The name stems from their characterization as a "fair weather waterspout on land". Waterspouts and landspouts share many defining characteristics, including relative weakness, short lifespan, and a small, smooth condensation funnel that often does not reach the surface. Landspouts also create a distinctively laminar cloud of dust when they make contact with the ground, due to their differing mechanics from true mesoform tornadoes. Though usually weaker than classic tornadoes, they can produce strong winds which could cause serious damage. Similar circulations Gustnado A gustnado, or gust front tornado, is a small, vertical swirl associated with a gust front or downburst. Because they are not connected with a cloud base, there is some debate as to whether or not gustnadoes are tornadoes. They are formed when fast-moving cold, dry outflow air from a thunderstorm is blown through a mass of stationary, warm, moist air near the outflow boundary, resulting in a "rolling" effect (often exemplified through a roll cloud). If low level wind shear is strong enough, the rotation can be turned vertically or diagonally and make contact with the ground. The result is a gustnado. They usually cause small areas of heavier rotational wind damage among areas of straight-line wind damage. Dust devil A dust devil (also known as a whirlwind) resembles a tornado in that it is a vertical swirling column of air. However, they form under clear skies and are no stronger than the weakest tornadoes. They form when a strong convective updraft is formed near the ground on a hot day. If there is enough low-level wind shear, the column of hot, rising air can develop a small cyclonic motion that can be seen near the ground. They are not considered tornadoes because they form during fair weather and are not associated with any clouds. However, they can, on occasion, result in major damage. Fire whirls Small-scale, tornado-like circulations can occur near any intense surface heat source. Those that occur near intense wildfires are called fire whirls. They are not considered tornadoes, except in the rare case where they connect to a pyrocumulus or other cumuliform cloud above. Fire whirls usually are not as strong as tornadoes associated with thunderstorms. They can, however, produce significant damage. Steam devils A steam devil is a rotating updraft between that involves steam or smoke. These formations do not involve high wind speeds, only completing a few rotations per minute. Steam devils are very rare. They most often form from smoke issuing from a power plant's smokestack. Hot springs and deserts may also be suitable locations for a tighter, faster-rotating steam devil to form. The phenomenon can occur over water, when cold arctic air passes over relatively warm water. Intensity and damage The Fujita scale, Enhanced Fujita scale (EF), and International Fujita scale rate tornadoes by damage caused. The EF scale was an update to the older Fujita scale, by expert elicitation, using engineered wind estimates and better damage descriptions. The EF scale was designed so that a tornado rated on the Fujita scale would receive the same numerical rating, and was implemented starting in the United States in 2007. An EF0 tornado will probably damage trees but not substantial structures, whereas an EF5 tornado can rip buildings off their foundations leaving them bare and even deform large skyscrapers. The similar TORRO scale ranges from a T0 for extremely weak tornadoes to T11 for the most powerful known tornadoes. Doppler weather radar data, photogrammetry, and ground swirl patterns (cycloidal marks) may also be analyzed to determine intensity and award a rating. Tornadoes vary in intensity regardless of shape, size, and location, though strong tornadoes are typically larger than weak tornadoes. The association with track length and duration also varies, although longer track tornadoes tend to be stronger. In the case of violent tornadoes, only a small portion of the path is of violent intensity, most of the higher intensity from subvortices. In the United States, 80% of tornadoes are EF0 and EF1 (T0 through T3) tornadoes. The rate of occurrence drops off quickly with increasing strength—less than 1% are violent tornadoes (EF4, T8 or stronger). Current records may significantly underestimate the frequency of strong (EF2-EF3) and violent (EF4-EF5) tornadoes, as damage-based intensity estimates are limited to structures and vegetation that a tornado impacts. A tornado may be much stronger than its damage-based rating indicates if its strongest winds occur away from suitable damage indicators, such as in an open field. Outside Tornado Alley, and North America in general, violent tornadoes are extremely rare. This is apparently mostly due to the lesser number of tornadoes overall, as research shows that tornado intensity distributions are fairly similar worldwide. A few significant tornadoes occur annually in Europe, Asia, southern Africa, and southeastern South America. Climatology The United States has the most tornadoes of any country, nearly four times more than estimated in all of Europe, excluding waterspouts. This is mostly due to the unique geography of the continent. North America is a large continent that extends from the tropics north into arctic areas, and has no major east–west mountain range to block air flow between these two areas. In the middle latitudes, where most tornadoes of the world occur, the Rocky Mountains block moisture and buckle the atmospheric flow, forcing drier air at mid-levels of the troposphere due to downsloped winds, and causing the formation of a low pressure area downwind to the east of the mountains. Increased westerly flow off the Rockies force the formation of a dry line when the flow aloft is strong, while the Gulf of Mexico fuels abundant low-level moisture in the southerly flow to its east. This unique topography allows for frequent collisions of warm and cold air, the conditions that breed strong, long-lived storms throughout the year. A large portion of these tornadoes form in an area of the central United States known as Tornado Alley. This area extends into Canada, particularly Ontario and the Prairie Provinces, although southeast Quebec, the interior of British Columbia, and western New Brunswick are also tornado-prone. Tornadoes also occur across northeastern Mexico. The United States averages about 1,200 tornadoes per year, followed by Canada, averaging 62 reported per year. NOAA's has a higher average 100 per year in Canada. The Netherlands has the highest average number of recorded tornadoes per area of any country (more than 20, or annually), followed by the UK (around 33, per year), although those are of lower intensity, briefer and cause minor damage. Tornadoes kill an average of 179 people per year in Bangladesh, the most in the world. Reasons for this include the region's high population density, poor construction quality, and lack of tornado safety knowledge. Other areas of the world that have frequent tornadoes include South Africa, the La Plata Basin area, portions of Europe, Australia and New Zealand, and far eastern Asia. Tornadoes are most common in spring and least common in winter, but tornadoes can occur any time of year that favorable conditions occur. Spring and fall experience peaks of activity as those are the seasons when stronger winds, wind shear, and atmospheric instability are present. Tornadoes are focused in the right front quadrant of landfalling tropical cyclones, which tend to occur in the late summer and autumn. Tornadoes can also be spawned as a result of eyewall mesovortices, which persist until landfall. Tornados can even form during snow squalls events with no rain present. Tornado occurrence is highly dependent on the time of day, because of solar heating. Worldwide, most tornadoes occur in the late afternoon, between 15:00 (3 pm) and 19:00 (7 pm) local time, with a peak near 17:00 (5 pm). Destructive tornadoes can occur at any time of day. The Gainesville Tornado of 1936, one of the deadliest tornadoes in history, occurred at 8:30 am local time. The United Kingdom has the highest incidence of tornadoes per unit area of land in the world. Unsettled conditions and weather fronts transverse the British Isles at all times of the years, and are responsible for spawning the tornadoes, which consequently form at all times of the year. The United Kingdom has at least 34 tornadoes per year and possibly as many as 50. Most tornadoes in the United Kingdom are weak, but they are occasionally destructive. For example, the Birmingham tornado of 2005 and the London tornado of 2006 both registered F2 on the Fujita scale and both caused significant damage and injury. Associations with climate and climate change Associations with various climate and environmental trends exist. For example, an increase in the sea surface temperature of a source region (e.g. Gulf of Mexico and Mediterranean Sea) increases atmospheric moisture content. Increased moisture can fuel an increase in severe weather and tornado activity, particularly in the cool season. Some evidence does suggest that the Southern Oscillation is weakly correlated with changes in tornado activity, which vary by season and region, as well as whether the ENSO phase is that of El Niño or La Niña. Research has found that fewer tornadoes and hailstorms occur in winter and spring in the U.S. central and southern plains during El Niño, and more occur during La Niña, than in years when temperatures in the Pacific are relatively stable. Ocean conditions could be used to forecast extreme spring storm events several months in advance. Climatic shifts may affect tornadoes via teleconnections in shifting the jet stream and the larger weather patterns. The climate-tornado link is confounded by the forces affecting larger patterns and by the local, nuanced nature of tornadoes. Although it is reasonable to suspect that global warming may affect trends in tornado activity, any such effect is not yet identifiable due to the complexity, local nature of the storms, and database quality issues. Any effect would vary by region. Detection Rigorous attempts to warn of tornadoes began in the United States in the mid-20th century. Before the 1950s, the only method of detecting a tornado was by someone seeing it on the ground. Often, news of a tornado would reach a local weather office after the storm. However, with the advent of weather radar, areas near a local office could get advance warning of severe weather. The first public tornado warnings were issued in 1950 and the first tornado watches and convective outlooks came about in 1952. In 1953, it was confirmed that hook echoes were associated with tornadoes. By recognizing these radar signatures, meteorologists could detect thunderstorms probably producing tornadoes from several miles away. Radar Today most developed countries have a network of weather radars, which serves as the primary method of detecting hook signatures that are likely associated with tornadoes. In the United States and a few other countries, Doppler weather radar stations are used. These devices measure the velocity and radial direction (towards or away from the radar) of the winds within a storm, and so can spot evidence of rotation in storms from over away. When storms are distant from a radar, only areas high within the storm are observed and the important areas below are not sampled. Data resolution also decreases with distance from the radar. Some meteorological situations leading to tornadogenesis are not readily detectable by radar and tornado development may occasionally take place more quickly than radar can complete a scan and send the batch of data. Doppler weather radar systems can detect mesocyclones within a supercell thunderstorm. This allows meteorologists to predict tornado formations throughout thunderstorms. Storm spotting Spotters usually are trained by the NWS on behalf of their respective organizations, and report to them. The organizations activate public warning systems such as sirens and the Emergency Alert System (EAS), and they forward the report to the NWS. There are more than 230,000 trained Skywarn weather spotters across the United States. In Canada, a similar network of volunteer weather watchers, called Canwarn, helps spot severe weather, with more than 1,000 volunteers. In Europe, several nations are organizing spotter networks under the auspices of Skywarn Europe and the Tornado and Storm Research Organisation (TORRO) has maintained a network of spotters in the United Kingdom since 1974. Storm spotters are required because radar systems such as NEXRAD detect signatures that suggest the presence of tornadoes, rather than tornadoes as such. Radar may give a warning before there is any visual evidence of a tornado or an imminent one, but ground truth from an observer can give definitive information. The spotter's ability to see what radar cannot is especially important as distance from the radar site increases, because the radar beam becomes progressively higher in altitude further away from the radar, chiefly due to curvature of Earth, and the beam also spreads out. Visual evidence Storm spotters are trained to discern whether or not a storm seen from a distance is a supercell. They typically look to its rear, the main region of updraft and inflow. Under that updraft is a rain-free base, and the next step of tornadogenesis is the formation of a rotating wall cloud. The vast majority of intense tornadoes occur with a wall cloud on the backside of a supercell. Evidence of a supercell is based on the storm's shape and structure, and cloud tower features such as a hard and vigorous updraft tower, a persistent, large overshooting top, a hard anvil (especially when backsheared against strong upper level winds), and a corkscrew look or striations. Under the storm and closer to where most tornadoes are found, evidence of a supercell and the likelihood of a tornado includes inflow bands (particularly when curved) such as a "beaver tail", and other clues such as strength of inflow, warmth and moistness of inflow air, how outflow- or inflow-dominant a storm appears, and how far is the front flank precipitation core from the wall cloud. Tornadogenesis is most likely at the interface of the updraft and rear flank downdraft, and requires a balance between the outflow and inflow. Only wall clouds that rotate spawn tornadoes, and they usually precede the tornado between five and thirty minutes. Rotating wall clouds may be a visual manifestation of a low-level mesocyclone. Barring a low-level boundary, tornadogenesis is highly unlikely unless a rear flank downdraft occurs, which is usually visibly evidenced by evaporation of cloud adjacent to a corner of a wall cloud. A tornado often occurs as this happens or shortly afterwards; first, a funnel cloud dips and in nearly all cases by the time it reaches halfway down, a surface swirl has already developed, signifying a tornado is on the ground before condensation connects the surface circulation to the storm. Tornadoes may also develop without wall clouds, under flanking lines and on the leading edge. Spotters watch all areas of a storm, and the cloud base and surface. Extremes The tornado which holds most records in history was the Tri-State Tornado, which roared through parts of Missouri, Illinois, and Indiana on March 18, 1925. It was likely an F5, though tornadoes were not ranked on any scale in that era. It holds records for longest path length (), longest duration (about 3.5 hours), and fastest forward speed for a significant tornado () anywhere on Earth. In addition, it is the deadliest single tornado in United States history (695 dead). The tornado was also the costliest tornado in history at the time (unadjusted for inflation), but in the years since has been surpassed by several others if population changes over time are not considered. When costs are normalized for wealth and inflation, it ranks third today. The deadliest tornado in world history was the Daultipur-Salturia Tornado in Bangladesh on April 26, 1989, which killed approximately 1,300 people. One of the most extensive tornado outbreaks on record was the 1974 Super Outbreak, which affected a large area of the central United States and extreme southern Ontario on April 3 and 4, 1974. The outbreak featured 148 tornadoes in 18 hours, many of which were violent; seven were of F5 intensity, and twenty-three peaked at F4 strength. Sixteen tornadoes were on the ground at the same time during its peak. More than 300 people, possibly as many as 330, were killed. While direct measurement of the most violent tornado wind speeds is nearly impossible, since conventional anemometers would be destroyed by the intense winds and flying debris, some tornadoes have been scanned by mobile Doppler radar units, which can provide a good estimate of the tornado's winds. The highest wind speed ever measured in a tornado, which is also the highest wind speed ever recorded on the planet, is 301 ± 20 mph (484 ± 32 km/h) in the F5 Bridge Creek-Moore, Oklahoma, tornado which killed 36 people. The reading was taken about above the ground. Storms that produce tornadoes can feature intense updrafts, sometimes exceeding . Debris from a tornado can be lofted into the parent storm and carried a very long distance. A tornado which affected Great Bend, Kansas, in November 1915, was an extreme case, where a "rain of debris" occurred from the town, a sack of flour was found away, and a cancelled check from the Great Bend bank was found in a field outside of Palmyra, Nebraska, to the northeast. Waterspouts and tornadoes have been advanced as an explanation for instances of raining fish and other animals. Safety Though tornadoes can strike in an instant, there are precautions and preventative measures that can be taken to increase the chances of survival. Authorities such as the Storm Prediction Center in the United States advise having a pre-determined plan should a tornado warning be issued. When a warning is issued, going to a basement or an interior first-floor room of a sturdy building greatly increases chances of survival. In tornado-prone areas, many buildings have underground storm cellars, which have saved thousands of lives. Some countries have meteorological agencies which distribute tornado forecasts and increase levels of alert of a possible tornado (such as tornado watches and warnings in the United States and Canada). Weather radios provide an alarm when a severe weather advisory is issued for the local area, mainly available only in the United States. Unless the tornado is far away and highly visible, meteorologists advise that drivers park their vehicles far to the side of the road (so as not to block emergency traffic), and find a sturdy shelter. If no sturdy shelter is nearby, getting low in a ditch is the next best option. Highway overpasses are one of the worst places to take shelter during tornadoes, as the constricted space can be subject to increased wind speed and funneling of debris underneath the overpass. Myths and misconceptions Folklore often identifies a green sky with tornadoes, and though the phenomenon may be associated with severe weather, there is no evidence linking it specifically with tornadoes. It is often thought that opening windows will lessen the damage caused by the tornado. While there is a large drop in atmospheric pressure inside a strong tornado, the pressure difference is unlikely to cause significant damage. Opening windows may instead increase the severity of the tornado's damage. A violent tornado can destroy a house whether its windows are open or closed. Another commonly held misconception is that highway overpasses provide adequate shelter from tornadoes. This belief is partly inspired by widely circulated video captured during the 1991 tornado outbreak near Andover, Kansas, where a news crew and several other people took shelter under an overpass on the Kansas Turnpike and safely rode out a tornado as it passed nearby. However, a highway overpass is a dangerous place during a tornado, and the subjects of the video remained safe due to an unlikely combination of events: the storm in question was a weak tornado, the tornado did not directly strike the overpass, and the overpass itself was of a unique design. Due to the Venturi effect, tornadic winds are accelerated in the confined space of an overpass. Indeed, in the 1999 Oklahoma tornado outbreak of May 3, 1999, three highway overpasses were directly struck by tornadoes, and at each of the three locations there was a fatality, along with many life-threatening injuries. By comparison, during the same tornado outbreak, more than 2,000 homes were completely destroyed and another 7,000 damaged, and yet only a few dozen people died in their homes. An old belief is that the southwest corner of a basement provides the most protection during a tornado. The safest place is the side or corner of an underground room opposite the tornado's direction of approach (usually the northeast corner), or the central-most room on the lowest floor. Taking shelter in a basement, under a staircase, or under a sturdy piece of furniture such as a workbench further increases the chances of survival. There are areas which people believe to be protected from tornadoes, whether by being in a city, near a major river, hill, or mountain, or even protected by supernatural forces. Tornadoes have been known to cross major rivers, climb mountains, affect valleys, and have damaged several city centers. As a general rule, no area is safe from tornadoes, though some areas are more susceptible than others. Ongoing research Meteorology is a relatively young science and the study of tornadoes is newer still. Although researched for about 140 years and intensively so for around 60 years, there are still aspects of tornadoes which remain a mystery. Meteorologists have a fairly good understanding of the development of thunderstorms and mesocyclones, and the meteorological conditions conducive to their formation. However, the step from supercell, or other respective formative processes, to tornadogenesis and the prediction of tornadic vs. non-tornadic mesocyclones is not yet well known and is the focus of much research. Also under study are the low-level mesocyclone and the stretching of low-level vorticity which tightens into a tornado, in particular, what are the processes and what is the relationship of the environment and the convective storm. Intense tornadoes have been observed forming simultaneously with a mesocyclone aloft (rather than succeeding mesocyclogenesis) and some intense tornadoes have occurred without a mid-level mesocyclone. In particular, the role of downdrafts, particularly the rear-flank downdraft, and the role of baroclinic boundaries, are intense areas of study. Reliably predicting tornado intensity and longevity remains a problem, as do details affecting characteristics of a tornado during its life cycle and tornadolysis. Other rich areas of research are tornadoes associated with mesovortices within linear thunderstorm structures and within tropical cyclones. Meteorologists still do not know the exact mechanisms by which most tornadoes form, and occasional tornadoes still strike without a tornado warning being issued. Analysis of observations including both stationary and mobile (surface and aerial) in-situ and remote sensing (passive and active) instruments generates new ideas and refines existing notions. Numerical modeling also provides new insights as observations and new discoveries are integrated into our physical understanding and then tested in computer simulations which validate new notions as well as produce entirely new theoretical findings, many of which are otherwise unattainable. Importantly, development of new observation technologies and installation of finer spatial and temporal resolution observation networks have aided increased understanding and better predictions. Research programs, including field projects such as the VORTEX projects (Verification of the Origins of Rotation in Tornadoes Experiment), deployment of TOTO (the TOtable Tornado Observatory), Doppler on Wheels (DOW), and dozens of other programs, hope to solve many questions that still plague meteorologists. Universities, government agencies such as the National Severe Storms Laboratory, private-sector meteorologists, and the National Center for Atmospheric Research are some of the organizations very active in research; with various sources of funding, both private and public, a chief entity being the National Science Foundation. The pace of research is partly constrained by the number of observations that can be taken; gaps in information about the wind, pressure, and moisture content throughout the local atmosphere; and the computing power available for simulation. Solar storms similar to tornadoes have been recorded, but it is unknown how closely related they are to their terrestrial counterparts.	Physical sciences	Earth science	null
37556	https://en.wikipedia.org/wiki/Asperger%20syndrome	Asperger syndrome	Asperger syndrome (AS), also known as Asperger's syndrome or Asperger's, is a diagnostic label that has been used to describe a neurodevelopmental disorder characterized by significant difficulties in social interaction and nonverbal communication, along with restricted, repetitive patterns of behavior and interests. Asperger syndrome has been merged with other conditions into autism spectrum disorder (ASD) and is no longer a diagnosis in the WHO's ICD-11 or the APA's DSM-5-TR. It was considered milder than other diagnoses which were merged into ASD due to relatively unimpaired spoken language and intelligence. The syndrome was named in 1976 by English psychiatrist Lorna Wing after the Austrian pediatrician Hans Asperger, who, in 1944, described children in his care who struggled to form friendships, did not understand others' gestures or feelings, engaged in one-sided conversations about their favorite interests, and were clumsy. In 1990 (coming into effect in 1993), the diagnosis of Asperger syndrome was included in the tenth edition (ICD-10) of the World Health Organization's International Classification of Diseases, and in 1994, it was also included in the fourth edition (DSM-4) of the American Psychiatric Association's Diagnostic and Statistical Manual of Mental Disorders. However, with the publication of DSM-5 in 2013 the syndrome was removed, and the symptoms are now included within autism spectrum disorder along with classic autism and pervasive developmental disorder not otherwise specified (PDD-NOS). It was similarly merged into autism spectrum disorder in the International Classification of Diseases (ICD-11) in 2018 (published, coming into effect in 2022). The exact cause of autism, including what was formerly known as Asperger syndrome, is not well understood. While it has high heritability, the underlying genetics have not been determined conclusively. Environmental factors are also believed to play a role. Brain imaging has not identified a common underlying condition. There is no single treatment, and the UK's National Health Service (NHS) guidelines suggest that "treatment" of any form of autism should not be a goal, since autism is not "a disease that can be removed or cured". According to the Royal College of Psychiatrists, while co-occurring conditions might require treatment, "management of autism itself is chiefly about the provision of the education, training, and social support/care required to improve the person's ability to function in the everyday world". The effectiveness of particular interventions for autism is supported by only limited data. Interventions may include social skills training, cognitive behavioral therapy, physical therapy, speech therapy, parent training, and medications for associated problems, such as mood or anxiety. Autistic characteristics tend to become less obvious in adulthood, but social and communication difficulties usually persist. In 2015, Asperger syndrome was estimated to affect 37.2 million people globally, or about 0.5% of the population. The exact percentage of people affected has still not been firmly established. Autism spectrum disorder is diagnosed in males more often than females, and females are typically diagnosed at a later age. The modern conception of Asperger syndrome came into existence in 1981 and went through a period of popularization. It became a standardized diagnosis in the 1990s and was merged into ASD in 2013. Many questions and controversies about the condition remain. Classification The extent of the overlap between Asperger syndrome and other forms of autism, particularly what was sometimes called high-functioning autism is unclear. The ASD classification is to some extent an artifact of how autism was discovered, and it may not reflect the true nature of the spectrum; methodological problems have beset Asperger syndrome as a valid diagnosis from the outset. As noted above, in the 2010s, Asperger syndrome, as a separate diagnosis, was eliminated and folded into autism spectrum disorder in the DSM-5 and the ICD-11. Like the diagnosis of Asperger syndrome, the change was controversial. The World Health Organization (WHO) previously defined Asperger syndrome (AS) as one of the pervasive developmental disorders (PDD), which are a spectrum of psychological disorders that are characterized by abnormalities of social interaction and communication that pervade the individual's functioning, and by restricted and repetitive interests and behavior. Like other neurodevelopmental conditions, ASD begins in infancy or childhood, has a steady course without remission or relapse, and has impairments that result from maturation-related changes in various systems of the brain. Characteristics As a pervasive developmental disorder, Asperger syndrome is distinguished by a pattern of symptoms rather than a single symptom. It is characterized by qualitative impairment in social interaction, by stereotyped and restricted patterns of behavior, activities, and interests, and by no clinically significant delay in cognitive development or general delay in language. Intense preoccupation with a narrow subject, one-sided verbosity, restricted prosody, and physical clumsiness are typical of the condition, but are not required for diagnosis. Suicidal thoughts and behaviors are a serious concern within the autistic population. One study found that adults with Asperger syndrome exhibited suicidal thoughts at 9 times the rate of the general population. Of autistic study participants, 66% had experienced suicidal ideation, while 35% had planned or attempted suicide. Social interaction A lack of demonstrated empathy affects aspects of social relatability for persons with Asperger syndrome. Individuals with Asperger syndrome experience difficulties in basic elements of social interaction, which may include a failure to develop friendships or to seek shared enjoyments or achievements with others (e.g., showing others objects of interest); a lack of social or emotional reciprocity; and impaired nonverbal behaviors in areas such as eye contact, facial expression, posture, and gesture. People with Asperger syndrome may not be as withdrawn around others, compared with those with other forms of autism; they approach others, even if awkwardly. For example, a person with Asperger syndrome may engage in a one-sided, long-winded speech about a favorite topic, while misunderstanding or not recognizing the listener's feelings or reactions, such as a wish to change the topic of talk or end the interaction. This social awkwardness has been called "active but odd". Such failures to react appropriately to social interaction may appear as disregard for other people's feelings and may come across as rude or insensitive. However, not all individuals with Asperger syndrome will approach others. Some may even display selective mutism, not speaking at all to most people and excessively to specific others. The cognitive ability of children with Asperger syndrome often allows them to articulate social norms in a laboratory context, where they may be able to show a theoretical understanding of other people's emotions; however, they typically have difficulty acting on this knowledge in fluid, real-life situations. People with Asperger syndrome may analyze and distill their observations of social interaction into rigid behavioral guidelines and apply these rules in awkward ways, such as forced eye contact, resulting in a demeanor that appears rigid or socially naïve. A history of failed attempts to establish reciprocal social relationships can cause autistic individuals to isolate themselves and cease attempts to engage; however, autistic people overwhelmingly report a desire for social contact and friendship. Violent or criminal behavior The hypothesis that individuals with Asperger syndrome are predisposed to violent or criminal behavior has been investigated but is not supported by data. More evidence suggests that children diagnosed with Asperger syndrome are more likely to be victims, rather than offenders. A 2008 review found that about 80% of reported violent criminals with Asperger syndrome also had other coexisting psychotic psychiatric disorders such as schizoaffective disorder. The sample size of this review was small (n = 37). Empathy People with an Asperger profile might not be recognized for their empathetic qualities, due to variation in the ways empathy is felt and expressed. Some people feel deep empathy, but do not outwardly communicate these sentiments through facial expressions or language. Some people come to empathy through intellectual processes, using logic and reasoning to arrive at the feelings. People with Asperger profiles may be bullied or excluded by peers, and might as a result be guarded around people, which could appear as lack of empathy. People with Asperger profiles can still be caring individuals; indeed, it is particularly common for those with the profile to feel and exhibit deep concern for individual rights, human welfare, animal rights, environmental protection, and other global and humanitarian causes. Evidence suggests that in the "double empathy problem model, autistic people have a unique interaction style which is significantly more readable by other autistic people, compared to non-autistic people." Restricted and repetitive interests and behavior People with Asperger syndrome can display behavior, interests, and activities that are restricted and repetitive and are sometimes abnormally intense or focused. They may stick to inflexible routines, move in stereotyped and repetitive ways, preoccupy themselves with parts of objects, or engage in compulsive behaviors like lining objects up to form patterns. The pursuit of specific and narrow areas of interest is one of the most striking among possible features of AS. Individuals with AS may collect volumes of detailed information on a relatively narrow topic such as weather data or star names without necessarily having a genuine understanding of the broader topic. For example, a child might memorize camera model numbers while caring little about photography. This behavior is usually apparent by age five or six. Although these special interests may change from time to time, they typically become more unusual and narrowly focused and often dominate social interaction so much that the entire family may become immersed. Because narrow topics often capture the interest of children, this symptom may go unrecognized. Stereotyped and repetitive motor behaviors, called stimming, are a core part of the diagnosis of AS and other ASDs. Stims are believed to be used for self-soothing and regulate sensory input. They include hand movements such as flapping or twisting, and complex whole-body movements. These are typically repeated in longer bursts and look more voluntary or ritualistic than tics, which are usually faster, less rhythmical, and less often symmetrical. Stimming may have a connection with tics, and studies have reported a consistent comorbidity between AS and Tourette syndrome in the range of 8–20%, with one figure as high as 80% for tics of some kind or another, for which several explanations have been put forward, including common genetic factors and dopamine, glutamate, or serotonin abnormalities. According to the Adult Asperger Assessment (AAA) diagnostic test, a lack of interest in fiction and a positive preference towards non-fiction is common among adults with AS. Speech and language Although individuals with Asperger syndrome acquire language skills without significant general delay and their speech typically lacks significant abnormalities, language acquisition and use is often atypical. Abnormalities include verbosity; abrupt transitions; literal interpretations and miscomprehension of nuance; use of metaphor meaningful only to the speaker; auditory perception deficits; unusually pedantic, formal, or idiosyncratic speech; and oddities in loudness, pitch, intonation, prosody, and rhythm. Echolalia has also been observed in individuals with AS. Three aspects of communication patterns are of clinical interest: poor prosody, tangential and circumstantial speech, and marked verbosity. Although inflection and intonation may be less rigid or monotonic than in classic autism, people with AS often have a limited range of intonation: speech may be unusually fast, jerky, or loud. Speech may convey a sense of incoherence; the conversational style often includes monologues about topics that bore the listener, fails to provide context for comments, or fails to suppress internal thoughts. Individuals with AS may fail to detect whether the listener is interested or engaged in the conversation. The speaker's conclusion or point may never be made, and attempts by the listener to elaborate on the speech's content or logic, or to shift to related topics, are often unsuccessful. Children with AS may have a sophisticated vocabulary at a young age and such children have often been colloquially called "little professors" but have difficulty understanding figurative language and tend to use language literally. Children with AS appear to have particular weaknesses in areas of nonliteral language that include humor, irony, teasing, and sarcasm. Although individuals with AS usually understand the cognitive basis of humor, they seem to lack understanding of the intent of humor to share the enjoyment with others. Despite strong evidence of impaired humor appreciation, anecdotal reports of humor in individuals with AS seem to challenge some psychological theories of AS and autism. Motor and sensory perception Individuals with Asperger syndrome may have signs or symptoms that are independent of the diagnosis but can affect the individual or the family. These include differences in perception and problems with motor skills, sleep, and emotions. Individuals with AS often have excellent auditory and visual perception. Children with ASD often demonstrate enhanced perception of small changes in patterns such as arrangements of objects or well-known images; typically this is domain-specific and involves processing of fine-grained features. Conversely, compared with individuals with high-functioning autism, individuals with AS have deficits in some tasks involving visual-spatial perception, auditory perception, or visual memory. Many accounts of individuals with AS and ASD report other unusual sensory and perceptual skills and experiences. They may be unusually sensitive or insensitive to sound, light, and other stimuli; these sensory responses are found in other developmental disorders and are not specific to AS or to ASD. There is little support for increased fight-or-flight response or failure of habituation in autism; there is more evidence of decreased responsiveness to sensory stimuli, although several studies show no differences. Hans Asperger's initial accounts and other diagnostic schemes include descriptions of physical clumsiness. Children with AS may be delayed in acquiring skills requiring dexterity, such as riding a bicycle or opening a jar, and may seem to move awkwardly or feel "uncomfortable in their own skin". They may be poorly coordinated or have an odd or bouncy gait or posture, poor handwriting, or problems with motor coordination. They may show problems with proprioception (sensation of body position) on measures of developmental coordination disorder (motor planning disorder), balance, tandem gait, and finger-thumb apposition. There is no evidence that these motor skills problems differentiate AS from other high-functioning ASDs. Children with AS are more likely to have sleep problems, including difficulty in falling asleep, frequent nocturnal awakenings, and early morning awakenings. AS is also associated with high levels of alexithymia, which is difficulty in identifying and describing one's emotions. Although AS, lower sleep quality, and alexithymia are associated with each other, their causal relationship is unclear. Causes Hans Asperger described common traits among his patients' family members, especially fathers, and research supports this observation and suggests a genetic contribution to Asperger syndrome. Although no specific genetic factor has yet been identified, multiple factors are believed to play a role in the expression of autism, given the variability in symptoms seen in children. Hundreds of genes have been linked to AS, and these genes play crucial role in a multitude of biological processes, exerting influence over the maturation and functioning of the brain. Evidence for a genetic link is that AS tends to run in families where more family members have limited behavioral symptoms similar to AS (for example, some problems with social interaction, or with language and reading skills). Most behavioral genetic research suggests that all autism spectrum disorders have shared genetic mechanisms. There may be shared genes in which particular alleles make an individual vulnerable, and varying combinations result in differing severity and symptoms in each person with AS. A few ASD cases have been linked to exposure to teratogens (agents that cause birth defects) during the first eight weeks from conception. Although this does not exclude the possibility that ASD can be initiated or affected later, it is strong evidence that ASD arises very early in development. Many environmental factors have been hypothesized to act after birth, but none has been confirmed by scientific investigation. These environmental elements can act as independent and significant risk factors, or they can potentially influence pre-existing genetic factors in people who have a genetic predisposition. Mechanism Asperger syndrome appears to result from developmental factors that affect many or all functional brain systems, as opposed to localized effects. Although the specific underpinnings of AS or factors that distinguish it from other ASDs are unknown, and no clear pathology common to individuals with AS has emerged, it is still possible that AS's mechanism is separate from other ASDs. Neuroanatomical studies and the associations with teratogens strongly suggest that the mechanism includes alteration of brain development soon after conception. Abnormal fetal development may affect the final structure and connectivity of the brain, resulting in altered neural circuits controlling thought and behavior. Several theories of mechanism are available; none are likely to provide a complete explanation. General-processing theories One general-processing theory is weak central coherence theory, which hypothesizes that a limited ability to see the big picture underlies the central disturbance in ASD. A related theory—enhanced perceptual functioning—focuses more on the superiority of locally oriented and perceptual operations in autistic individuals. Mirror neuron system (MNS) theory The mirror neuron system (MNS) theory hypothesizes that alterations to the development of the MNS interfere with imitation and lead to Asperger syndrome's core feature of social impairment. One study found that activation is delayed in the core circuit for imitation in individuals with AS. This theory maps well to social cognition theories like the theory of mind, which hypothesizes that autistic behavior arises from impairments in ascribing mental states to oneself and others; or hyper-systemizing, which hypothesizes that autistic individuals can systematize internal operation to handle internal events but are less effective at empathizing when handling events generated by other agents. Diagnosis Standard diagnostic criteria require impairment in social interaction and repetitive and stereotyped patterns of behavior, activities, and interests, without significant delay in language or cognitive development. Unlike the international standard, the DSM-IV-TR criteria also required significant impairment in day-to-day functioning; As noted above, in the 2010s, Asperger syndrome, as a separate diagnosis, was eliminated and folded into autism spectrum disorder in the DSM-5 and the ICD-11. Other sets of diagnostic criteria have been proposed by Szatmari et al. and by Gillberg and Gillberg. Diagnosis of ASD (and previously AS) is most commonly made between the ages of four and eleven. A comprehensive assessment involves a multidisciplinary team that observes across multiple settings, and includes neurological and genetic assessment as well as tests for cognition, psychomotor function, verbal and nonverbal strengths and weaknesses, style of learning, and skills for independent living. The "gold standard" in diagnosing ASDs combines clinical judgment with the Autism Diagnostic Interview-Revised (ADI-R), a semistructured parent interview; and the Autism Diagnostic Observation Schedule (ADOS), a conversation and play-based interview with the child. Delayed or mistaken diagnosis can be traumatic for individuals and families; for example, misdiagnosis can lead to medications that worsen behavior. Underdiagnosis and overdiagnosis may be problems. The cost and difficulty of screening and assessment can delay diagnosis. Conversely, the increasing popularity of drug treatment options and the expansion of benefits has motivated providers to overdiagnose ASD. There are indications AS has been diagnosed more frequently in recent years, partly as a residual diagnosis for children of normal intelligence who are not autistic but have social difficulties. There are questions about the external validity of the AS diagnosis. That is, it is unclear whether there is a practical benefit in distinguishing AS from autism or PDD-NOS; different screening tools may render different diagnoses for the same person. Differential diagnosis Many children with AS are initially misdiagnosed with attention deficit hyperactivity disorder (ADHD). Diagnosing adults is more challenging, as standard diagnostic criteria are designed for children and the expression of AS changes with age. Adult diagnosis requires painstaking clinical examination and thorough medical history gained from both the individual and other people who know the person, focusing on childhood behavior. Conditions that must be considered in a differential diagnosis along with ADHD include other ASDs, the schizophrenia spectrum, personality disorders, obsessive–compulsive disorder, major depressive disorder, semantic pragmatic disorder, nonverbal learning disorder, social anxiety disorder, Tourette syndrome, stereotypic movement disorder, bipolar disorder, social-cognitive deficits due to brain damage from alcohol use disorder, and obsessive–compulsive personality disorder (OCPD). Screening Parents of children with Asperger syndrome can typically trace differences in their children's development to as early as 30 months of age. Developmental screening during a routine check-up by a general practitioner or pediatrician may identify signs that warrant further investigation. The United States Preventive Services Task Force in 2016 found it was unclear if screening was beneficial or harmful among children in whom there are no concerns. Different screening instruments are used to diagnose AS, including the Asperger Syndrome Diagnostic Scale (ASDS); Autism Spectrum Screening Questionnaire (ASSQ); Childhood Autism Spectrum Test (CAST), previously called the Childhood Asperger Syndrome Test; Gilliam Asperger's disorder scale (GADS); Krug Asperger's Disorder Index (KADI); and the autism-spectrum quotient (AQ), with versions for children, adolescents, and adults. None have been shown to reliably differentiate between AS and other ASDs. Management Treatment attempts to manage distressing symptoms and to teach age-appropriate social, communication, and vocational skills that are not naturally acquired during development. Intervention is tailored to the needs of the individual based on multidisciplinary assessment. Although progress has been made, data supporting the efficacy of particular interventions are limited. Therapies Managing ASD may involve multiple therapies that address core symptoms of the disorder. While many professionals agree that the earlier the professional support the better, there is no combination that is recommended above others. Professional support for ASD varies depending on the individual; it takes into account the linguistic capabilities, verbal strengths, and nonverbal vulnerabilities of individuals. Many of those diagnosed with ASD or similar disorders advocate against behavioral therapies, like Applied behavior analysis (ABA) and Cognitive behavioral therapy (CBT), often as part of the autism rights movement, on the grounds that these approaches frequently reinforce the demand on autistic people to mask their neurodivergent characteristics or behaviors to favor a more 'neurotypical' and narrow conception of normality. ABA has faced a great deal of criticism over the years. Recently, studies have shown that ABA may be abusive and can increase PTSD symptoms in patients. The Autistic Self Advocacy Network campaigns against the use of ABA in autism. In the case of CBT and talking therapies, the effectiveness varies, with many reporting that they appeared 'too self-aware' to gain significant benefit, as the therapy was designed with neurotypical people in mind. In autistic children, specifically, they also report that it is only mildly beneficial in aiding with their anxieties. A typical program of professional support generally includes: Applied behavior analysis (ABA) procedures, including positive behavior support (PBS)—or training and support of parents and school faculty in behavior management strategies to use in the home and school, and social skills training for more effective interpersonal interactions. The Autistic Self Advocacy Network campaigns against the use of ABA in autism; Cognitive behavioral therapy to improve stress management relating to anxiety or explosive emotions and to help reduce obsessive interests (although this may produce negative impact by demonising special interests) and repetitive routines; Medication for coexisting conditions such as major depressive disorder and anxiety disorders; Occupational or physical therapy to assist with poor sensory processing and motor coordination; and, Social communication intervention, which is specialized speech therapy to help with the pragmatics and give-and-take of normal conversation. Of the many studies on behavior-based early intervention programs, most are case reports of up to five participants and typically examine a few problem behaviors such as self-injury, aggression, noncompliance, stereotypies, or spontaneous language; unintended side effects are largely ignored. Despite the popularity of social skills training, its effectiveness is not firmly established. A randomized controlled study of a model for training parents in problem behaviors in their children with AS showed that parents attending a one-day workshop or six individual lessons reported fewer behavioral problems, while parents receiving the individual lessons reported less intense behavioral problems in their AS children. Vocational training may be important to teach job interview etiquette and workplace behavior to older children and adults with AS, and organization software and personal data assistants can improve the work and life management of people with AS. Fecal Microbiota Transplantation (FMT) is an innovative therapy for AS that aims to restore microbial balance in the patient's gastrointestinal tract by introducing healthy fecal microbiota acquired from people with a diverse microbial composition. This approach attempts to reconstruct the patient's gut microbiota by taking into account the intricate interactions between the human gut and the central nervous system via the gut-brain axis (GBA). Any disruption in gut health has been linked to an increased susceptibility to diverse neurodevelopmental disorders. It is vital to remember that research of AS specifically operates upon the out-dated classification of this syndrome as external to ASD (Autism Spectrum Disorder). Similarly, we should also note that ASD is a spectrum and support varies dramatically depending on the individual. Medications No medications directly treat the core symptoms of AS. Although research into the efficacy of pharmaceutical intervention for AS is limited, it is essential to diagnose and treat comorbid conditions. Deficits in self-identifying emotions or in observing effects of one's behavior on others can make it difficult for individuals with AS to see why medication may be appropriate. Medication can be effective in combination with behavioral interventions and environmental accommodations in treating comorbid symptoms such as anxiety disorders, major depressive disorder, inattention, and aggression. The atypical antipsychotic medications risperidone, olanzapine and aripiprazole have been shown to reduce the associated symptoms of AS; risperidone can reduce repetitive and self-injurious behaviors, aggressive outbursts, and impulsivity, and improve stereotypical patterns of behavior and social relatedness. The selective serotonin reuptake inhibitors (SSRIs) fluoxetine, fluvoxamine, and sertraline have been effective in treating restricted and repetitive interests and behaviors, while stimulant medication, such as methylphenidate, can reduce inattention. In addition, scientists have made a noteworthy finding that oxytocin, a hormone, plays a significant role in shaping human social behavior and the formation of interpersonal connections. Care must be taken with medications, as side effects may be more common and harder to evaluate in individuals with AS, and tests of drugs' effectiveness against comorbid conditions routinely exclude individuals from the autism spectrum. Abnormalities in metabolism, cardiac conduction times, and an increased risk of type 2 diabetes have been raised as concerns with antipsychotic medications, along with serious long-term neurological side effects. SSRIs can lead to manifestations of behavioral activation such as increased impulsivity, aggression, and sleep disturbance. Weight gain and fatigue are commonly reported side effects of risperidone, which may also lead to increased risk for extrapyramidal symptoms such as restlessness and dystonia and increased serum prolactin levels. Sedation and weight gain are more common with olanzapine, which has also been linked with diabetes. Sedative side-effects in school-age children have ramifications for classroom learning. Individuals with AS may be unable to identify and communicate their internal moods and emotions or to tolerate side effects that for most people would not be problematic. Prognosis There is some evidence that children with AS may see a lessening of symptoms; up to 20% of children may no longer meet the diagnostic criteria as adults, although social and communication difficulties may persist. , no studies addressing the long-term outcome of individuals with Asperger syndrome are available and there are no systematic long-term follow-up studies of children with AS. Individuals with AS appear to have normal life expectancy, but have an increased prevalence of comorbid psychiatric conditions, such as major depressive disorder and anxiety disorders that may significantly affect prognosis. Although social impairment may be lifelong, the outcome is generally more positive than with individuals with lower-functioning autism spectrum disorders; for example, ASD symptoms are more likely to diminish with time in children with AS or forms of autism sometimes described as "high functioning". Most students with AS and forms of autism sometimes seen as "high functioning" have average mathematical ability and test slightly worse in mathematics than in general intelligence. However, mathematicians are at least three times more likely to have autism-spectrum traits than the general population, and are more likely to have family members with autism. Although many attend regular education classes, some children with AS may attend special education classes such as separate classroom and resource room because of their social and behavioral difficulties. Adolescents with AS may exhibit ongoing difficulty with self-care or organization, and disturbances in social and romantic relationships. Despite high cognitive potential, most young adults with AS remain at home, yet some do marry and work independently. The "different-ness" adolescents experience can be traumatic. Anxiety may stem from preoccupation over possible violations of routines and rituals, from being placed in a situation without a clear schedule or expectations, or from concern with failing in social encounters; the resulting stress may manifest as inattention, withdrawal, reliance on obsessions, hyperactivity, or aggressive or oppositional behavior. Depression is often the result of chronic frustration from repeated failure to engage others socially, and mood disorders requiring treatment may develop. Clinical experience suggests the rate of suicide may be higher among those with AS, but this has not been confirmed by systematic empirical studies. Education of families is critical in developing strategies for understanding strengths and weaknesses; helping the family to cope improves outcomes in children. Prognosis may be improved by diagnosis at a younger age that allows for early interventions, while interventions in adulthood are valuable but less beneficial. There are legal implications for individuals with AS as they run the risk of exploitation by others and may be unable to comprehend the societal implications of their actions. Epidemiology Frequency estimates vary enormously. In 2015, it was estimated that 37.2 million people globally are affected. A 2003 review of epidemiological studies of children found autism rates ranging from 0.03 to 4.84 per 1,000, with the ratio of autism to Asperger syndrome ranging from 1.5:1 to 16:1; combining the geometric mean ratio of 5:1 with a conservative prevalence estimate for autism of 1.3 per 1,000 suggests indirectly that the prevalence of AS might be around 0.26 per 1,000. Part of the variance in estimates arises from differences in diagnostic criteria. For example, a relatively small 2007 study of 5,484 eight-year-old children in Finland found 2.9 children per 1,000 met the ICD-10 criteria for an AS diagnosis, 2.7 per 1,000 for Gillberg and Gillberg criteria, 2.5 for DSM-IV, 1.6 for Szatmari et al., and 4.3 per 1,000 for the union of the four criteria. Boys seem to be more likely to have AS than girls; estimates of the sex ratio range from 1.6:1 to 4:1, using the Gillberg and Gillberg criteria. Females with autism spectrum disorders may be underdiagnosed. Comorbidities Anxiety disorders and major depressive disorder are the most common conditions seen at the same time; comorbidity of these in persons with AS is estimated at 65%. Reports have associated AS with medical conditions such as aminoaciduria and ligamentous laxity, but these have been case reports or small studies and no factors have been associated with AS across studies. One study of males with AS found an increased rate of epilepsy and a high rate (51%) of nonverbal learning disorder. AS is associated with tics, Tourette syndrome and bipolar disorder. The repetitive behaviors of AS have many similarities with the symptoms of obsessive–compulsive disorder and obsessive–compulsive personality disorder, and 26% of a sample of young adults with AS were found to meet the criteria for schizoid personality disorder (which is characterised by severe social seclusion and emotional detachment), more than any other personality disorder in the sample. However many of these studies are based on clinical samples or lack standardized measures; nonetheless, comorbid conditions are relatively common. Correlated characteristics Research indicates that individuals with Aspergers have significantly higher rates of LGBT identities and feelings than the general population. They are also significantly more likely to be non-theistic. History Asperger syndrome was named after the Austrian pediatrician Hans Asperger (1906–1980), but not coined by him. Asperger syndrome was a relatively new diagnosis in the field of autism, though a syndrome like it was described as early as 1925 by Soviet child psychiatrist Grunya Sukhareva (1891–1981), As a child, Asperger appears to have exhibited some features of the very condition named after him, such as remoteness and talent in language. In 1944, Asperger gave detailed descriptions of four representative children in his practice who had difficulty in integrating themselves socially and showing empathy towards peers. They also lacked nonverbal communication skills and were physically clumsy. Asperger described this "autistic psychopathy" as social isolation. Fifty years later, several standardizations of AS as a medical diagnosis were tentatively proposed, many of which diverge significantly from Asperger's original work. Unlike what became known as AS, Asperger believed autistic psychopathy could be found in people of all levels of intelligence, including those with intellectual disability: as such, Asperger's understanding of autistic pathology was more akin to what is known as the autism spectrum today. Asperger defended the value of so-called "high-functioning" autistic individuals, writing: "We are convinced, then, that autistic people have their place in the organism of the social community. They fulfill their role well, perhaps better than anyone else could, and we are talking of people who as children had the greatest difficulties and caused untold worries to their care-givers." Asperger also believed some would be capable of exceptional achievement and original thought later in life. Asperger's paper was published during World War II and in German, so it was not widely read elsewhere. Lorna Wing used the term Asperger syndrome in 1976, and popularized it to the English-speaking medical community in her February 1981 publication of case studies of children showing the symptoms described by Asperger, and Uta Frith translated Asperger's paper to English in 1991. Sets of diagnostic criteria were outlined by Gillberg and Gillberg in 1989 and by Szatmari et al. in the same year. AS became a standard diagnosis when it was included in the tenth edition of the World Health Organization's diagnostic manual, International Classification of Diseases (ICD-10), published in 1990 and coming into effect in 1993; and in the fourth edition of the American Psychiatric Association's diagnostic reference, Diagnostic and Statistical Manual of Mental Disorders (DSM-IV), published in 1994. Hundreds of books, articles, and websites later described AS and prevalence estimates increased dramatically for ASD, with AS recognized as an important subgroup. Whether AS should be seen as distinct from autism, particularly forms of autism sometimes described as sometimes described as "high functioning", became an issue receiving significant attention and disagreement, along with questions about the empirical validation of the DSM-IV and ICD-10 criteria. With the publication of the next major editions of the DSM and ICD, the DSM-5 (published in 2013) and the ICD-11 (published in 2018, coming into effect in 2022), AS was eliminated as a separate diagnosis and folded into the autism spectrum. A scale of "severity" levels was included in the DSM-5, whereby most people previously diagnosed with AS would have been classified as "level 1"; but these levels are widely opposed by the autistic community and are not included in the ICD-11. The ICD-11 characterizes ASD with qualifiers describing the presence of disorders of intellectual development and the degree of functional language impairment; the former diagnosis of Asperger syndrome is characterized as autism spectrum disorder without disorder of intellectual development and with mild or no impairment of functional language. Society and culture People identifying with Asperger syndrome may refer to themselves in casual conversation as aspies (a term first used in print in the Boston Globe in 1998). Some autistic people have advocated a shift in perception of autism spectrum disorders as complex syndromes, neurodivergences, and/or neurominority cognitive styles rather than diseases that must be cured. Proponents of this neurodiversity paradigm reject the notion that there is an "ideal" brain configuration and that any deviation from the norm is pathological; they promote tolerance of neurodiversity. These views are the basis for the autistic rights and autistic pride movements, within the broader neurodiversity movement. There is a contrast between the attitude of people with AS, who typically do not want to be cured and are proud of their identity; and parents of children with AS, who more often seek a "cure" of their children's autism. Some researchers have argued that AS and other autism can be viewed as a different cognitive style, not a disorder, and that it should be removed from psychiatric and medical manuals classifying diseases (ICD) or mental disorders (DSM), much as homosexuality was removed. Even some people typically associated with a pathology paradigm for autism are willing to consider AS a neutral difference. For example, in 2002, Simon Baron-Cohen wrote of those with AS: "In the social world, there is no great benefit to a precise eye for detail, but in the worlds of maths, computing, cataloging, music, linguistics, engineering, and science, such an eye for detail can lead to success rather than failure." Baron-Cohen cited two reasons why it might still be useful to consider AS to be a disability: to ensure provision for legally required special support, and to recognize emotional difficulties from reduced empathy, which was commonly associated with autism during that time but has since lost support. Baron-Cohen argues that the genes for ASD's combination of abilities have operated throughout recent human evolution and have made remarkable contributions to human history. By contrast, Pier Jaarsma and Welin wrote in 2011 that the "broad version of the neurodiversity claim, covering low-functioning as well as high-functioning autism, is problematic. Only a narrow conception of neurodiversity, referring exclusively to high-functioning autists, is reasonable." They say that "higher functioning" individuals with autism may "not [be] benefited with such a psychiatric defect-based diagnosis ... some of them are being harmed by it, because of the disrespect the diagnosis displays for their natural way of being", but "think that it is still reasonable to include other categories of autism in the psychiatric diagnostics. The narrow conception of the neurodiversity claim should be accepted but the broader claim should not."	Biology and health sciences	Mental disorders	Health
37575	https://en.wikipedia.org/wiki/Airport	Airport	An airport is an aerodrome with extended facilities, mostly for commercial air transport. They usually consist of a landing area, which comprises an aerially accessible open space including at least one operationally active surface such as a runway for a plane to take off and to land or a helipad, and often includes adjacent utility buildings such as control towers, hangars and terminals, to maintain and monitor aircraft. Larger airports may have airport aprons, taxiway bridges, air traffic control centres, passenger facilities such as restaurants and lounges, and emergency services. In some countries, the US in particular, airports also typically have one or more fixed-base operators, serving general aviation. Airport operations are extremely complex, with a complicated system of aircraft support services, passenger services, and aircraft control services contained within the operation. Thus airports can be major employers, as well as important hubs for tourism and other kinds of transit. Because they are sites of operation for heavy machinery, a number of regulations and safety measures have been implemented in airports, in order to reduce hazards. Additionally, airports have major local environmental impacts, as both large sources of air pollution, noise pollution and other environmental impacts, making them sites that acutely experience the environmental effects of aviation. Airports are also vulnerable infrastructure to extreme weather, climate change caused sea level rise and other disasters. Terminology The terms aerodrome, airfield, and airstrip also refer to airports, and the terms heliport, seaplane base, and STOLport refer to airports dedicated exclusively to helicopters, seaplanes, and short take-off and landing aircraft. In colloquial use in certain environments, the terms airport and aerodrome are often interchanged. However, in general, the term airport may imply or confer a certain stature upon the aviation facility that other aerodromes may not have achieved. In some jurisdictions, airport is a legal term of art reserved exclusively for those aerodromes certified or licensed as airports by the relevant civil aviation authority after meeting specified certification criteria or regulatory requirements. That is to say, all airports are aerodromes, but not all aerodromes are airports. In jurisdictions where there is no legal distinction between aerodrome and airport, which term to use in the name of an aerodrome may be a commercial decision. In US technical/legal usage, landing area is used instead of aerodrome, and airport means "a landing area used regularly by aircraft for receiving or discharging passengers or cargo". Types of airports An airport solely serving helicopters is called a heliport. An airport for use by seaplanes and amphibious aircraft is called a seaplane base. Such a base typically includes a stretch of open water for takeoffs and landings, and seaplane docks for tying-up. An international airport has additional facilities for customs and passport control as well as incorporating all the aforementioned elements. Such airports rank among the most complex and largest of all built typologies, with 15 of the top 50 buildings by floor area being airport terminals. Management Smaller or less-developed airfields, which represent the vast majority, often have a single runway shorter than . Larger airports for airline flights generally have paved runways of or longer. Skyline Airport in Inkom, Idaho, has a runway that is only long. In the United States, the minimum dimensions for dry, hard landing fields are defined by the FAR Landing And Takeoff Field Lengths. These include considerations for safety margins during landing and takeoff. The longest public-use runway in the world is at Qamdo Bamda Airport in China. It has a length of . The world's widest paved runway is at Ulyanovsk Vostochny Airport in Russia and is wide. , the CIA stated that there were approximately 44,000 "airports or airfields recognizable from the air" around the world, including 15,095 in the US, the US having the most in the world. Airport ownership and operation Most of the world's large airports are owned by local, regional, or national government bodies who then lease the airport to private corporations who oversee the airport's operation. For example, in the UK the state-owned British Airports Authority originally operated eight of the nation's major commercial airports – it was subsequently privatized in the late 1980s, and following its takeover by the Spanish Ferrovial consortium in 2006, has been further divested and downsized to operating just Heathrow. Germany's Frankfurt Airport is managed by the quasi-private firm Fraport. While in India GMR Group operates, through joint ventures, Indira Gandhi International Airport and Rajiv Gandhi International Airport. Bengaluru International Airport is controlled by Fairfax .Chhatrapati Shivaji International Airport, Chaudhary Charan Singh International Airport, Mangalore International Airport, Thiruvananthapuram International Airport, Lokpriya Gopinath Bordoloi International Airport, Jaipur International Airport, Sardar Vallabhbhai Patel International Airport are operated by Adani Group through a Public Private Partnership wherein Adani Group, the operator pays Airports Authority of India, the owner of the airports, a predetermined sum of money based on the number of passengers handled by the airports. The rest of India's airports are managed by the Airports Authority of India. In Pakistan nearly all civilian airports are owned and operated by the Pakistan Civil Aviation Authority except for Sialkot International Airport which has the distinction of being the first privately owned public airport in Pakistan and South Asia. In the US, commercial airports are generally operated directly by government entities or government-created airport authorities (also known as port authorities), such as the Los Angeles World Airports authority that oversees several airports in the Greater Los Angeles area, including Los Angeles International Airport. In Canada, the federal authority, Transport Canada, divested itself of all but the remotest airports in 1999/2000. Now most airports in Canada are operated by individual legal authorities, such as Vancouver International Airport Authority (although still owned by Transport Canada); some airports, such as Boundary Bay Airport and Pitt Meadows Airport, are municipally owned. Many US airports still lease part or all of their facilities to outside firms, who operate functions such as retail management and parking. All US commercial airport runways are certified by the FAA under the Code of Federal Regulations Title 14 Part 139, "Certification of Commercial Service Airports" but maintained by the local airport under the regulatory authority of the FAA. Despite the reluctance to privatize airports in the US (contrary to the FAA sponsoring a privatization program since 1996), the government-owned, contractor-operated (GOCO) arrangement is the standard for the operation of commercial airports in the rest of the world. Airport funding The Airport & Airway Trust Fund (AATF) was created by the Airport and Airway Development in 1970 which finances aviation programs in the United States. Airport Improvement Program (AIP), Facilities and Equipment (F&E), and Research, Engineering, and Development (RE&D) are the three major accounts of Federal Aviation Administration which are financed by the AATF, as well as pays for the FAA's Operation and Maintenance (O&M) account. The funding of these accounts are dependent on the taxes the airports generate of revenues. Passenger tickets, fuel, and cargo tax are the taxes that are paid by the passengers and airlines help fund these accounts. Airport revenue Airports revenues are divided into three major parts: aeronautical revenue, non-aeronautical revenue, and non-operating revenue. Aeronautical revenue makes up 50% in 2021 (from 54% and 48% in 2019 and 2020, non-aeronautical revenue makes up 34% (40%, 39% in previous years), and non-operating revenue makes up 16% (6%, 14%) of the total revenue of airports. Aeronautical revenue Aeronautical revenue are generated through airline rents and landing, passenger service, parking, and hangar fees. Landing fees are charged per aircraft for landing an airplane in the airport property. Landing fees are calculated through the landing weight and the size of the aircraft which varies but most of the airports have a fixed rate and a charge extra for extra weight. Passenger service fees are charges per passengers for the facilities used on a flight like water, food, wifi and shows which is paid while paying for an airline ticket. Aircraft parking is also a major revenue source for airports. Aircraft are parked for a certain amount of time before or after takeoff and have to pay to park there. Every airport has its own rates of parking, for example, John F Kennedy airport in New York City charges $45 per hour for a plane of 100,000 pounds and the price increases with weight. Non-aeronautical revenue Non-aeronautical revenue is gained through things other than aircraft operations. It includes lease revenue from compatible land-use development, non-aeronautical building leases, retail and concession sales, rental car operations, parking and in-airport advertising. Concession revenue is one big part of non-aeronautical revenue airports makes through duty free, bookstores, restaurants and money exchange. Car parking is a growing source of revenue for airports, as more people use the parking facilities of the airport. O'Hare International Airport in Chicago charges $2 per hour for every car. Price regulation Many airports are local monopolies. To prevent them from abusing their market power, governments regulate how much airports may charge to airlines, using price-cap regulation. Landside and airside areas Airports are divided into landside and airside zones. The landside is subject to fewer special laws and is part of the public realm, while access to the airside zone is tightly controlled. Landside facilities may include publicly accessible airport check-in desks, shops and ground transportation facilities. The airside area includes all parts of the airport around the aircraft, and the parts of the buildings that are restricted to staff, and sections of these extended to travelling, airside shopping, dining, or waiting passengers. Depending on the airport, passengers and staff must be checked by security or border control before being permitted to enter the airside zone. Conversely, passengers arriving from an international flight must pass through border control and customs to access the landside area, in which they exit, unless in airside transit. Most multi-terminal airports have (variously termed) flight/passenger/air connections buses, moving walkways and/or people movers for inter-terminal airside transit. Their airlines can arrange for baggage to be routed directly to the passenger's destination. Most major airports issue a secure keycard, an airside pass to employees, to assist in their reliable, standardized and efficient verification of identity. Facilities A terminal is a building with passenger facilities. Small airports have one terminal. Large ones often have multiple terminals, though some large airports, like Amsterdam Airport Schiphol, still have one terminal. The terminal has a series of gates, which provide passengers with access to the plane. Passenger facilities typically include: Airport check-in, including a baggage drop-off Airport security Passport control Gates Waiting areas Baggage reclaim facilities, often in the form of a carousel Customs (international arrivals only) Links between passenger facilities and aircraft include jet bridges or airstairs. Baggage handling systems transport baggage from the baggage drop-off to departing planes, and from arriving planes to the baggage reclaim. The area where the aircraft parks to load passengers and baggage is known as an apron or ramp (or incorrectly, "the tarmac"). Airport security Airport security normally requires baggage checks, metal screenings of individual persons, and rules against any object that could be used as a weapon. Since the September 11 attacks and the Real ID Act of 2005, airport security has dramatically increased and gotten tighter and stricter than ever before. Products and services Most major airports provide commercial outlets for products and services. Most of these companies, many of which are internationally known brands, are located within the departure areas. These include clothing boutiques and restaurants and in the US amounted to $4.2 billion in 2015. Prices charged for items sold at these outlets are generally higher than those outside the airport. However, some airports now regulate costs to keep them comparable to "street prices". This term is misleading as prices often match the manufacturers' suggested retail price (MSRP) but are almost never discounted. Many new airports include walkthrough duty-free stores that require air passengers to enter a retail store upon exiting security. Airport planners sometimes incorporate winding routes within these stores such that passengers encounter more goods as they walk towards their gate. Planners also install artworks next to the airport's shops in order to draw passengers into the stores. Apart from major fast food chains, some airport restaurants offer regional cuisine specialties for those in transit so that they may sample local food without leaving the airport. Some airport structures include on-site hotels built within or attached to a terminal building. Airport hotels have grown popular due to their convenience for transient passengers and easy accessibility to the airport terminal. Many airport hotels also have agreements with airlines to provide overnight lodging for displaced passengers. Major airports in such countries as Russia and Japan offer miniature sleeping units within the airport that are available for rent by the hour. The smallest type is the capsule hotel popular in Japan. A slightly larger variety is known as a sleep box. An even larger type is provided by the company YOTEL. Some airports provide smoking areas and prayer areas. Premium and VIP services Airports may also contain premium and VIP services. The premium and VIP services may include express check-in and dedicated check-in counters. These services are usually reserved for first and business class passengers, premium frequent flyers, and members of the airline's clubs. Premium services may sometimes be open to passengers who are members of a different airline's frequent flyer program. This can sometimes be part of a reciprocal deal, as when multiple airlines are part of the same alliance, or as a ploy to attract premium customers away from rival airlines. Sometimes these premium services will be offered to a non-premium passenger if the airline has made a mistake in handling of the passenger, such as unreasonable delays or mishandling of checked baggage. Airline lounges frequently offer free or reduced cost food, as well as alcoholic and non-alcoholic beverages. Lounges themselves typically have seating, showers, quiet areas, televisions, computer, Wi-Fi and Internet access, and power outlets that passengers may use for their electronic equipment. Some airline lounges employ baristas, bartenders and gourmet chefs. Airlines sometimes operate multiple lounges within the one airport terminal allowing ultra-premium customers, such as first class customers, additional services, which are not available to other premium customers. Multiple lounges may also prevent overcrowding of the lounge facilities. Cargo and freight service In addition to people, airports move cargo around the clock. Cargo airlines often have their own on-site and adjacent infrastructure to transfer parcels between ground and air. Cargo Terminal Facilities are areas where international airports export cargo has to be stored after customs clearance and prior to loading the aircraft. Similarly, import cargo that is offloaded needs to be in bond before the consignee decides to take delivery. Areas have to be kept aside for examination of export and import cargo by the airport authorities. Designated areas or sheds may be given to airlines or freight forward ring agencies. Every cargo terminal has a landside and an airside. The landside is where the exporters and importers through either their agents or by themselves deliver or collect shipments while the airside is where loads are moved to or from the aircraft. In addition, cargo terminals are divided into distinct areas – export, import, and interline or transshipment. Access and onward travel Airports require parking lots, for passengers who may leave the cars at the airport for a long period of time. Large airports will also have car-rental firms, taxi ranks, bus stops and sometimes a train station. Many large airports are located near railway trunk routes for seamless connection of multimodal transport, for instance Frankfurt Airport, Amsterdam Airport Schiphol, London Heathrow Airport, Tokyo Haneda Airport, Tokyo Narita Airport, Hamad International Airport, London Gatwick Airport and London Stansted Airport. It is also common to connect an airport and a city with rapid transit, light rail lines or other non-road public transport systems. Some examples of this would include the AirTrain JFK at John F. Kennedy International Airport in New York, Link light rail that runs from the heart of downtown Seattle to Seattle–Tacoma International Airport, and the Silver Line T at Boston's Logan International Airport by the Massachusetts Bay Transportation Authority (MBTA). Such a connection lowers risk of missed flights due to traffic congestion. Large airports usually have access also through controlled-access highways ('freeways' or 'motorways') from which motor vehicles enter either the departure loop or the arrival loop. Internal transport The distances passengers need to move within a large airport can be substantial. It is common for airports to provide moving walkways, buses, and rail transport systems. Some airports like Hartsfield–Jackson Atlanta International Airport and London Stansted Airport have a transit system that connects some of the gates to a main terminal. Airports with more than one terminal have a transit system to connect the terminals together, such as John F. Kennedy International Airport, Mexico City International Airport and London Gatwick Airport. Airport operations Airport operations are made possible by an organized network of trained personnel, specialized equipment, and spatial data. After thousands of ground operations staff left the industry during the COVID-19 pandemic, there have been discussions on the need for systemic improvements in three primary areas: Improving talent acquisition and retention Implementing global standardizations Digitizing and automating processes The surfaces where ground operations occur are generally divided into three regions: runways, taxiways, and aprons. Air traffic control Air traffic control (ATC) is the task of managing aircraft movements and making sure they are safe, orderly and expeditious. At the largest airports, air traffic control is a series of highly complex operations that requires managing frequent traffic that moves in all three dimensions. A "towered" or "controlled" airport has a control tower where the air traffic controllers are based. Pilots are required to maintain two-way radio communication with the controllers, and to acknowledge and comply with their instructions. A "non-towered" airport has no operating control tower and therefore two-way radio communications are not required, though it is good operating practice for pilots to transmit their intentions on the airport's common traffic advisory frequency (CTAF) for the benefit of other aircraft in the area. The CTAF may be a Universal Integrated Community (UNICOM), MULTICOM, Flight Service Station (FSS), or tower frequency. The majority of the world's airports are small facilities without a tower. Not all towered airports have 24/7 ATC operations. In those cases, non-towered procedures apply when the tower is not in use, such as at night. Non-towered airports come under area (en-route) control. Remote and virtual tower (RVT) is a system in which ATC is handled by controllers who are not present at the airport itself. Air traffic control responsibilities at airports are usually divided into at least two main areas: ground and tower, though a single controller may work both stations. The busiest airports may subdivide responsibilities further, with clearance delivery, apron control, and/or other specialized ATC stations. Ground control Ground control is responsible for directing all ground traffic in designated "movement areas", except the traffic on runways. This includes planes, baggage trains, snowplows, grass cutters, fuel trucks, stair trucks, airline food trucks, conveyor belt vehicles and other vehicles. Ground Control will instruct these vehicles on which taxiways to use, which runway they will use (in the case of planes), where they will park, and when it is safe to cross runways. When a plane is ready to takeoff it will be turned over to tower control. Conversely, after a plane has landed it will depart the runway and be "handed over" from Tower to Ground Control. Tower control Tower control is responsible for aircraft on the runway and in the controlled airspace immediately surrounding the airport. Tower controllers may use radar to locate an aircraft's position in 3D space, or they may rely on pilot position reports and visual observation. They coordinate the sequencing of aircraft in the traffic pattern and direct aircraft on how to safely join and leave the circuit. Aircraft which are only passing through the airspace must also contact tower control to be sure they remain clear of other traffic. Traffic pattern At all airports the use of a traffic pattern (often called a traffic circuit outside the US) is possible. They may help to assure smooth traffic flow between departing and arriving aircraft. There is no technical need within modern commercial aviation for performing this pattern, provided there is no queue. And due to the so-called SLOT-times, the overall traffic planning tend to assure landing queues are avoided. If for instance an aircraft approaches runway 17 (which has a heading of approx. 170 degrees) from the north (coming from 360/0 degrees heading towards 180 degrees), the aircraft will land as fast as possible by just turning 10 degrees and follow the glidepath, without orbit the runway for visual reasons, whenever this is possible. For smaller piston engined airplanes at smaller airfields without ILS equipment, things are very different though. Generally, this pattern is a circuit consisting of five "legs" that form a rectangle (two legs and the runway form one side, with the remaining legs forming three more sides). Each leg is named (see diagram), and ATC directs pilots on how to join and leave the circuit. Traffic patterns are flown at one specific altitude, usually above ground level (AGL). Standard traffic patterns are left-handed, meaning all turns are made to the left. One of the main reason for this is that pilots sit on the left side of the airplane, and a Left-hand patterns improves their visibility of the airport and pattern. Right-handed patterns do exist, usually because of obstacles such as a mountain, or to reduce noise for local residents. The predetermined circuit helps traffic flow smoothly because all pilots know what to expect, and helps reduce the chance of a mid-air collision. At controlled airports, a circuit can be in place but is not normally used. Rather, aircraft (usually only commercial with long routes) request approach clearance while they are still hours away from the airport; the destination airport can then plan a queue of arrivals, and planes will be guided into one queue per active runway for a "straight-in" approach. While this system keeps the airspace free and is simpler for pilots, it requires detailed knowledge of how aircraft are planning to use the airport ahead of time and is therefore only possible with large commercial airliners on pre-scheduled flights. The system has recently become so advanced that controllers can predict whether an aircraft will be delayed on landing before it even takes off; that aircraft can then be delayed on the ground, rather than wasting expensive fuel waiting in the air. Navigational aids There are a number of aids, both visual and electronic, though not at all airports. A visual approach slope indicator (VASI) helps pilots fly the approach for landing. Some airports are equipped with a VHF omnidirectional range (VOR) to help pilots find the direction to the airport. VORs are often accompanied by a distance measuring equipment (DME) to determine the distance to the VOR. VORs are also located off airports, where they serve to provide airways for aircraft to navigate upon. In poor weather, pilots will use an instrument landing system (ILS) to find the runway and fly the correct approach, even if they cannot see the ground. The number of instrument approaches based on the use of the Global Positioning System (GPS) is rapidly increasing and may eventually become the primary means for instrument landings. Larger airports sometimes offer precision approach radar (PAR), but these systems are more common at military air bases than civilian airports. The aircraft's horizontal and vertical movement is tracked via radar, and the controller tells the pilot his position relative to the approach slope. Once the pilots can see the runway lights, they may continue with a visual landing. Taxiway signs Airport guidance signs provide direction and information to taxiing aircraft and airport vehicles. Smaller aerodromes may have few or no signs, relying instead on diagrams and charts. Lighting Many airports have lighting that help guide planes using the runways and taxiways at night or in rain or fog. On runways, green lights indicate the beginning of the runway for landing, while red lights indicate the end of the runway. Runway edge lighting consists of white lights spaced out on both sides of the runway, indicating the edges. Some airports have more complicated lighting on the runways including lights that run down the centerline of the runway and lights that help indicate the approach (an approach lighting system, or ALS). Low-traffic airports may use pilot-controlled lighting to save electricity and staffing costs. Along taxiways, blue lights indicate the taxiway's edge, and some airports have embedded green lights that indicate the centerline. Weather observations Weather observations at the airport are crucial to safe takeoffs and landings. In the United States and Canada, the vast majority of airports, large and small, will either have some form of automated airport weather station, whether an AWOS, ASOS, or AWSS, a human observer or a combination of the two. These weather observations, predominantly in the METAR format, are available over the radio, through automatic terminal information service (ATIS), via the ATC or the flight service station. Planes take-off and land into the wind to achieve maximum performance. Because pilots need instantaneous information during landing, a windsock can also be kept in view of the runway. Aviation windsocks are made with lightweight material, withstand strong winds and some are lit up after dark or in foggy weather. Because visibility of windsocks is limited, often multiple glow-orange windsocks are placed on both sides of the runway. Airport ground crew (ground handling) Most airports have groundcrew handling the loading and unloading of passengers, crew, baggage and other services. Some groundcrew are linked to specific airlines. Among the vehicles that serve an airliner on the ground are: A tow tractor to move the aircraft in and out of the berth. A jet bridge (in some airports) or stairs unit to allow passengers to embark and disembark. A ground power unit for supplying electricity. As the engines will be switched off, they will not be generating electricity as they do in flight. A cleaning service. A catering service to deliver food and drinks for a flight. A toilet waste truck to empty the tank which holds the waste from the toilets in the aircraft. A water truck to fill the water tanks of the aircraft. A refueling vehicle. The fuel may come from a tanker, or from underground fuel tanks. A conveyor belt unit for loading and unloading luggage. A vehicle to transport luggage to and from the terminal. The length of time an aircraft remains on the ground in between consecutive flights is known as "turnaround time". Airlines seek to minimize turnaround times, with times scheduled as low as 25 minutes. Maintenance management Like industrial equipment or facility management, airports require tailor-made maintenance management due to their complexity. With many tangible assets spread over a large area in different environments, these infrastructures must therefore effectively monitor these assets and store spare parts to maintain them at an optimal level of service. To manage these airport assets, several solutions are competing for the market: CMMS (computerized maintenance management system) predominate, and mainly enable a company's maintenance activity to be monitored, planned, recorded and rationalized. Safety management Aviation safety is an important concern in the operation of an airport, and almost every airfield includes equipment and procedures for handling emergency situations. Airport crash tender crews are equipped for dealing with airfield accidents, crew and passenger extractions, and the hazards of highly flammable aviation fuel. The crews are also trained to deal with situations such as bomb threats, hijacking, and terrorist activities. Hazards to aircraft include debris, nesting birds, and reduced friction levels due to environmental conditions such as ice, snow, or rain. Part of runway maintenance is airfield rubber removal which helps maintain friction levels. The fields must be kept clear of debris using cleaning equipment so that loose material does not become a projectile and enter an engine duct (see foreign object damage). In adverse weather conditions, ice and snow clearing equipment can be used to improve traction on the landing strip. For waiting aircraft, equipment is used to spray special deicing fluids on the wings. Many airports are built near open fields or wetlands. These tend to attract bird populations, which can pose a hazard to aircraft in the form of bird strikes. Airport crews often need to discourage birds from taking up residence. Some airports are located next to parks, golf courses, or other low-density uses of land. Other airports are located near densely populated urban or suburban areas. An airport can have areas where collisions between aircraft on the ground tend to occur. Records are kept of any incursions where aircraft or vehicles are in an inappropriate location, allowing these "hot spots" to be identified. These locations then undergo special attention by transportation authorities (such as the FAA in the US) and airport administrators. During the 1980s, a phenomenon known as microburst became a growing concern due to aircraft accidents caused by microburst wind shear, such as Delta Air Lines Flight 191. Microburst radar was developed as an aid to safety during landing, giving two to five minutes' warning to aircraft in the vicinity of the field of a microburst event. Some airfields now have a special surface known as soft concrete at the end of the runway (stopway or blastpad) that behaves somewhat like styrofoam, bringing the plane to a relatively rapid halt as the material disintegrates. These surfaces are useful when the runway is located next to a body of water or other hazard, and prevent the planes from overrunning the end of the field. Airports often have on-site firefighters to respond to emergencies. These use specialized vehicles, known as airport crash tenders. Most civil aviation authorities have required levels of on-site emergency response capabilities based on an airport's traffic. At airports where civil and military operations share a common set of runways and infrastructure, emergency response is often managed by the relevant military unit as part of their base's operations. Environmental concerns and sustainability Aircraft noise is a major cause of noise disturbance to residents living near airports. Sleep can be affected if the airports operate night and early morning flights. Aircraft noise occurs not only from take-offs and landings but also from ground operations including maintenance and testing of aircraft. Noise can have other health effects as well. Other noises and environmental concerns are vehicle traffic causing noise and pollution on roads leading to the airport. The construction of new airports or addition of runways to existing airports, is often resisted by local residents because of the effect on countryside, historical sites, and local flora and fauna. Due to the risk of collision between birds and aircraft, large airports undertake population control programs where they frighten or shoot birds. The construction of airports has been known to change local weather patterns. For example, because they often flatten out large areas, they can be susceptible to fog in areas where fog rarely forms. In addition, they generally replace trees and grass with pavement, they often change drainage patterns in agricultural areas, leading to more flooding, run-off and erosion in the surrounding land. Airports are often built on low-lying coastal land, globally 269 airports are at risk of coastal flooding now. A temperature rise of 2oC – consistent with the Paris Agreement – would lead to 100 airports being below mean sea level and 364 airports at risk of flooding. If global mean temperature rise exceeds this then as many as 572 airports will be at risk by 2100, leading to major disruptions without appropriate adaptation. Some of the airport administrations prepare and publish annual environmental reports to show how they consider these environmental concerns in airport management issues and how they protect environment from airport operations. These reports contain all environmental protection measures performed by airport administration in terms of water, air, soil and noise pollution, resource conservation and protection of natural life around the airport. A 2019 report from the Cooperative Research Programs of the US Transportation Research Board showed all airports have a role to play in advancing greenhouse gas (GHG) reduction initiatives. Small airports have demonstrated leadership by using their less complex organizational structure to implement newer technologies and to serve as a proving ground for their feasibility. Large airports have the economic stability and staff resources necessary to grow in-house expertise and fund comprehensive new programs. A growing number of airports are installing solar photovoltaic arrays to offset their electricity use. The National Renewable Energy Lab has shown this can be done safely. This can also be done on the roofs of the airports and it has been found that the solar panels on these buildings work more effectively when compared to residential panels. The world's first airport to be fully powered by solar energy is located at Kochi, India. Another airport known for considering environmental concerns is Seymour Airport in the Galapagos Islands. As a part of their sustainability efforts, more and more airports are starting to explore the consequences of more electric aircraft coming into service. Electric aircraft require much energy; operating 49 small 50-passenger short-range battery electric aircraft would demand at least 16 GWh/year, and with short turnaround times between different flights, the charging powers have to be substantial. To tackle these issues, more airports are starting to look into alternative energy production such as solar power and wind power, but also how to use airport areas for biomass production. Another solution investigated is to use energy storage to charge during the night and use to charge the aircraft during daytime. Airport hygiene and public health concerns Airports, as major international travel hubs, have the potential to be significant transmission points for infectious diseases. A notable study conducted during the peak of the 2015–2016 flu season at Helsinki-Vantaa airport in Finland revealed that commonly touched surfaces in airports, especially the plastic security screening trays, are highly susceptible to contamination by respiratory viruses. These trays are touched by hundreds of passengers daily and, being made of plastic, a non-porous material, provide an environment where viruses can survive for extended periods. In comparison, bathroom surfaces in the same airport tested negative for respiratory viruses, possibly indicating a heightened awareness of hygiene in these spaces. Hand hygiene plays a pivotal role in preventing the spread of infectious diseases in airports. Research indicates that only about 70% of individuals wash their hands after using the toilet, and of those, only 50% do so correctly. In airport settings, just one in five individuals maintain clean hands—defined as washing with soap for at least 15 seconds in the preceding hour. Given the frequent touching of shared surfaces in airports, such as trays, railings, and touch panels, this poses a significant risk for disease transmission. A study from the University of Cyprus and MIT used models and simulations to demonstrate that increasing hand cleanliness from 20% to 30% in all airports could reduce the potential global impact of a disease by 24%. If cleanliness levels reached 60%, this impact could decrease by 69%. Focusing on just the top 10 most influential airports for disease spread, enhancing hand hygiene practices could still significantly reduce disease transmission rates from 45% to 37%. The findings underscore the importance of promoting hand-washing in airports as a key measure in preventing the global spread of diseases. Military air base An airbase, sometimes referred to as an air station or airfield, provides basing and support of military aircraft. Some airbases, known as military airports, provide facilities similar to their civilian counterparts. For example, RAF Brize Norton in the UK has a terminal that caters to passengers for the Royal Air Force's scheduled flights to the Falkland Islands. Some airbases are co-located with civilian airports, sharing the same ATC facilities, runways, taxiways and emergency services, but with separate terminals, parking areas and hangars. Bardufoss Airport, Bardufoss Air Station in Norway and Pune Airport in India are examples of this. An aircraft carrier is a warship that functions as a mobile airbase. Aircraft carriers allow a naval force to project air power without having to depend on local bases for land-based aircraft. After their development in World War I, aircraft carriers replaced the battleship as the centrepiece of a modern fleet during World War II. Airport designation and naming Most airports in the United States are designated "private-use airports" meaning that, whether publicly- or privately owned, the airport is not open or available for use by the public (although use of the airport may be made available by invitation of the owner or manager). Most airport names include the location. Many airport names honour a public figure, commonly a politician (e.g., Charles de Gaulle Airport, George Bush Intercontinental Airport, Lennart Meri Airport, O.R. Tambo International Airport, Soekarno–Hatta International Airport), a monarch (e.g. Chhatrapati Shivaji International Airport, King Shaka International Airport), a cultural leader (e.g. Liverpool John Lennon Airport, Leonardo da Vinci-Fiumicino Airport, Louis Armstrong New Orleans International Airport) or a prominent figure in aviation history of the region (e.g. Sydney Kingsford Smith Airport), sometimes even famous writers (e.g. Allama Iqbal International Airport) and explorers (e.g. Venice Marco Polo Airport). Some airports have unofficial names, possibly so widely circulated that its official name is little used or even known. Some airport names include the word "International" to indicate their ability to handle international air traffic. This includes some airports that do not have scheduled international airline services (e.g. Port Elizabeth International Airport). History and development The earliest aircraft takeoff and landing sites were grassy fields. The plane could approach at any angle that provided a favorable wind direction. A slight improvement was the dirt-only field, which eliminated the drag from grass. However, these functioned well only in dry conditions. Later, concrete surfaces would allow landings regardless of meteorological conditions. The title of "world's oldest airport" is disputed. Toussus-le-Noble airport near Paris was established in 1907 and has been operating since. College Park Airport in Maryland, US, established in 1909 by Wilbur Wright serves only general aviation traffic. Beijing Nanyuan Airport in China, which was built to accommodate planes in 1904, and airships in 1907, opened in 1910. It was in operation until September 2019. Pearson Field Airport in Vancouver, Washington, United States, was built to accommodate planes in 1905 and airships in 1911, and is still in use as of February 2024. Hamburg Airport opened in January 1911, making it the oldest commercial airport in the world which is still in operation. Bremen Airport opened in 1913 and remains in use, although it served as an American military field between 1945 and 1949. Amsterdam Airport Schiphol opened on September 16, 1916, as a military airfield, but has accepted civil aircraft only since December 17, 1920, allowing Sydney Airport—which started operations in January 1920—to claim to be one of the world's oldest continuously operating commercial airports. Minneapolis-Saint Paul International Airport in the US opened in 1920 and has been in continuous commercial service since. It serves about 35,000,000 passengers each year and continues to expand, recently opening a new 11,000-foot (3,355 m) runway. Of the airports constructed during this early period in aviation, it is one of the largest and busiest that is still currently operating. Don Mueang International Airport near Bangkok, Thailand, opened 1914, is also a contender, as well as the Rome Ciampino Airport, which opened in 1916. Increased aircraft traffic during World War I led to the construction of landing fields. Aircraft had to approach these from certain directions and this led to the development of aids for directing the approach and landing slope. Following the war, some of these military airfields added civil facilities for handling passenger traffic. One of the earliest such fields was Paris – Le Bourget Airport at Le Bourget, near Paris. The first airport to operate scheduled international commercial services was Hounslow Heath Aerodrome in August 1919, but it was closed and supplanted by Croydon Airport in March 1920. In 1922, the first permanent airport and commercial terminal solely for commercial aviation was opened at Flughafen Devau near what was then Königsberg, East Prussia. The airports of this era used a paved "apron", which permitted night flying as well as landing heavier aircraft. The first lighting used on an airport was during the latter part of the 1920s; in the 1930s approach lighting came into use. These indicated the proper direction and angle of descent. The colours and flash intervals of these lights became standardized under the International Civil Aviation Organization (ICAO). In the 1940s, the slope-line approach system was introduced. This consisted of two rows of lights that formed a funnel indicating an aircraft's position on the glideslope. Additional lights indicated incorrect altitude and direction. After World War II, airport design became more sophisticated. Passenger buildings were being grouped together in an island, with runways arranged in groups about the terminal. This arrangement permitted expansion of the facilities. But it also meant that passengers had to travel further to reach their plane. An improvement in the landing field was the introduction of grooves in the concrete surface. These run perpendicular to the direction of the landing aircraft and serve to draw off excess rainwater that could build up in front of the plane's wheels. Airport construction boomed during the 1960s with the increase in jet aircraft traffic. Runways were extended out to . The fields were constructed out of reinforced concrete using a slip-form machine that produces a continuous slab with no disruptions along the length. The early 1960s also saw the introduction of jet bridge systems to modern airport terminals, an innovation which eliminated outdoor passenger boarding. These systems became commonplace in the United States by the 1970s. The malicious use of UAVs has led to the deployment of counter unmanned air system (C-UAS) technologies such as the Aaronia AARTOS which have been installed on major international airports. Airports in entertainment Airports have played major roles in films and television programs due to their very nature as a transport and international hub, and sometimes because of distinctive architectural features of particular airports. One such example of this is The Terminal, a film about a man who becomes permanently grounded in an airport terminal and must survive only on the food and shelter provided by the airport. They are also one of the major elements in movies such as The V.I.P.s, Speed, Airplane!, Airport (1970), Die Hard 2, Soul Plane, Jackie Brown, Get Shorty, Home Alone (1990), Home Alone 2: Lost in New York (1992), Liar Liar, Passenger 57, Final Destination (2000), Unaccompanied Minors, Catch Me If You Can, Rendition and The Langoliers. They have also played important parts in television series like Lost, The Amazing Race, America's Next Top Model (season 10), 90 Day Fiancé, Air Crash Investigation which have significant parts of their story set within airports. In other programmes and films, airports are merely indicative of journeys, e.g. Good Will Hunting. Several computer simulation games put the player in charge of an airport. These include the Airport Tycoon series, SimAirport and Airport CEO. Airport directories Each civil aviation authority provides a source of information about airports in their country. This will contain information on airport elevation, airport lighting, runway information, communications facilities and frequencies, hours of operation, nearby NAVAIDs and contact information where prior arrangement for landing is necessary. Australia Information can be found on-line in the En route Supplement Australia (ERSA) which is published by Airservices Australia, a government owned corporation charged with managing Australian ATC. Brazil Infraero is responsible for the airports in Brazil Canada Two publications, the Canada Flight Supplement (CFS) and the Water Aerodrome Supplement, published by Nav Canada under the authority of Transport Canada provides equivalent information. Europe The European Organisation for the Safety of Air Navigation (EUROCONTROL) provides an Aeronautical Information Publication (AIP), aeronautical charts and NOTAM services for multiple European countries. Germany Provided by the Luftfahrt-Bundesamt (Federal Office for Civil Aviation of Germany). France Aviation Generale Delage edited by Delville and published by Breitling. The United Kingdom The information is found in Pooley's Flight Guide, a publication compiled with the assistance of the United Kingdom Civil Aviation Authority (CAA). Pooley's also contains information on some continental European airports that are close to Great Britain. National Air Traffic Services, the UK's Air Navigation Service Provider, a public–private partnership also publishes an online AIP for the UK. The United States The US uses the Airport/Facility Directory (A/FD) (now officially termed the Chart Supplement) published in seven volumes. DAFIF also includes extensive airport data but has been unavailable to the public at large since 2006. Japan Aeronautical Information Publication (AIP) is provided by Japan Aeronautical Information Service Center, under the authority of Japan Civil Aviation Bureau, Ministry of Land, Infrastructure, Transport and Tourism of Japan. A comprehensive, consumer/business directory of commercial airports in the world (primarily for airports as businesses, rather than for pilots) is organized by the trade group Airports Council International.	Technology	Aviation	null
37583	https://en.wikipedia.org/wiki/Wine%20%28software%29	Wine (software)	Wine is a free and open-source compatibility layer to allow application software and computer games developed for Microsoft Windows to run on Unix-like operating systems. Developers can compile Windows applications against WineLib to help port them to Unix-like systems. Wine is predominantly written using black-box testing reverse-engineering, to avoid copyright issues. No code emulation or virtualization occurs. Wine is primarily developed for Linux and macOS. In a 2007 survey by desktoplinux.com of 38,500 Linux desktop users, 31.5% of respondents reported using Wine to run Windows applications. This plurality was larger than all x86 virtualization programs combined, and larger than the 27.9% who reported not running Windows applications. History Bob Amstadt, the initial project leader, and Eric Youngdale started the Wine project in 1993 as a way to run Windows applications on Linux. It was inspired by two Sun Microsystems products, Wabi for the Solaris operating system, and the Public Windows Interface, which was an attempt to get the Windows API fully reimplemented in the public domain as an ISO standard but rejected due to pressure from Microsoft in 1996. Wine originally targeted 16-bit applications for Windows 3.x, but focuses on 32-bit and 64-bit versions which have become the standard on newer operating systems. The project originated in discussions on Usenet in comp.os.linux in June 1993. Alexandre Julliard has led the project since 1994. The project has proven time-consuming and difficult for the developers, mostly because of incomplete and incorrect documentation of the Windows API. While Microsoft extensively documents most Win32 functions, some areas such as file formats and protocols have no public, complete specification available from Microsoft. Windows also includes undocumented low-level functions, undocumented behavior and obscure bugs that Wine must duplicate precisely in order to allow some applications to work properly. Consequently, the Wine team has reverse-engineered many function calls and file formats in such areas as thunking. The Wine project originally released Wine under the same MIT License as the X Window System, but owing to concern about proprietary versions of Wine not contributing their changes back to the core project, work as of March 2002 has used the LGPL for its licensing. Wine officially entered beta with version 0.9 on 25 October 2005. Version 1.0 was released on 17 June 2008, after 15 years of development. Version 1.2 was released on 16 July 2010, version 1.4 on 7 March 2012, version 1.6 on 18 July 2013, version 1.8 on 19 December 2015 and version 9.0 on 16 January 2024. Development versions are released roughly every two weeks. Wine-staging is an independently maintained set of aggressive patches not deemed ready by WineHQ developers for merging into the Wine repository, but still considered useful by the wine-compholio fork. It mainly covers experimental functions and bug fixes. Since January 2017, patches in wine-staging begins to be actively merged into the WineHQ upstream as wine-compholio transferred the project to Alistair Leslie-Hughes, a key WineHQ developer. , WineHQ also provides pre-built versions of wine-staging. Corporate sponsorship The main corporate sponsor of Wine is CodeWeavers, which employs Julliard and many other Wine developers to work on Wine and on CrossOver, CodeWeavers' supported version of Wine. CrossOver includes some application-specific tweaks not considered suitable for the upstream version, as well as some additional proprietary components. The involvement of Corel for a time assisted the project, chiefly by employing Julliard and others to work on it. Corel had an interest in porting WordPerfect Office, its office suite, to Linux (especially Corel Linux). Corel later cancelled all Linux-related projects after Microsoft made major investments in Corel, stopping their Wine effort. Other corporate sponsors include Google, which hired CodeWeavers to fix Wine so Picasa ran well enough to be ported directly to Linux using the same binary as on Windows; Google later paid for improvements to Wine's support for Adobe Photoshop CS2. Wine is also a regular beneficiary of Google's Summer of Code program. Valve works with CodeWeavers to develop Proton, a Wine-based compatibility layer for Microsoft Windows games to run on Linux-based operating systems. Proton includes several patches that upstream Wine does not accept for various reasons, such as Linux-specific implementations of Win32 functions. Valve's involvement in the development of Proton (and, thus, the improvement of Linux gaming) has helped to improve Wine compatibility with Windows games. Design The goal of Wine is to implement the Windows APIs fully or partially that are required by programs that the users of Wine wish to run on top of a Unix-like system. Basic architecture The programming interface of Microsoft Windows consists largely of dynamic-link libraries (DLLs). These contain a huge number of wrapper sub-routines for the system calls of the kernel, the NTOS kernel-mode program (ntoskrnl.exe). A typical Windows program calls some Windows DLLs, which in turn calls user-mode gdi/user32 libraries, which in turn uses the kernel32.dll (win32 subsystem) responsible for dealing with the kernel through system calls. The system-call layer is considered private to Microsoft programmers as documentation is not publicly available, and published interfaces all rely on subsystems running on top of the kernel. Besides these, there are a number of programming interfaces implemented as services that run as separate processes. Applications communicate with user-mode services through RPCs. Wine implements the Windows application binary interface (ABI) entirely in user space, rather than as a kernel module. Wine mostly mirrors the hierarchy, with services normally provided by the kernel in Windows instead provided by a daemon known as the wineserver, whose task is to implement basic Windows functionality, as well as integration with the X Window System, and translation of signals into native Windows exceptions. Although wineserver implements some aspects of the Windows kernel, it is not possible to use native Windows drivers with it, due to Wine's underlying architecture. Libraries and applications Wine allows for loading both Windows DLLs and Unix shared objects for its Windows programs. Its built-in implementation of the most basic Windows DLLs, namely NTDLL, KERNEL32, GDI32, and USER32, uses the shared object method because they must use functions in the host operating system as well. Higher-level libraries, such as WineD3D, are free to use the DLL format. In many cases users can choose to load a DLL from Windows instead of the one implemented by Wine. Doing so can provide functionalities not yet implemented by Wine, but may also cause malfunctions if it relies on something else not present in Wine. Wine tracks its state of implementation through automated unit testing done at every git commit. Graphics and gaming While most office software does not make use of complex GPU-accelerated graphics APIs, computer games do. To run these games properly, Wine would have to forward the drawing instructions to the host OS, and even translate them to something the host can understand. DirectX is a collection of Microsoft APIs for rendering, audio and input. As of 2019, Wine 4.0 contains a DirectX 12 implementation for Vulkan API, and DirectX 11.2 for OpenGL. Direct2D support has been updated to Direct2D 1.2. Wine 4.0 also allows Wine to run Vulkan applications by handing draw commands to the host OS, or in the case of macOS, by translating them into the Metal API by MoltenVK. XAudio , Wine 4.3 uses the FAudio library (and Wine 4.13 included a fix for it) to implement the XAudio2 audio API (and more). XInput and Raw Input Wine, since 4.0 (2019), supports game controllers through its builtin implementations of these libraries. They are built as Unix shared objects as they need to access the controller interfaces of the underlying OS, specifically through SDL. Direct3D Much of Wine's DirectX effort goes into building WineD3D, a translation layer from Direct3D and DirectDraw API calls into OpenGL. As of 2019, this component supports up to DirectX 11. As of 12 December 2016, Wine is good enough to run Overwatch with D3D11. Besides being used in Wine, WineD3D DLLs have also been used on Windows itself, allowing for older GPUs to run games using newer DirectX versions and for old DDraw-based games to render correctly. Some work is ongoing to move the Direct3D backend to Vulkan API. Direct3D 12 support in 4.0 is provided by a "vkd3d" subproject, and WineD3D has in 2019 been experimentally ported to use the Vulkan API. Another implementation, DXVK, translates Direct3D 8, 9, 10, and 11 calls using Vulkan as well and is a separate project. Wine, when patched, can alternatively run Direct3D 9 API commands directly via a free and open-source Gallium3D State Tracker (aka Gallium3D GPU driver) without translation into OpenGL API calls. In this case, the Gallium3D layer allows a direct pass-through of DX9 drawing commands which results in performance improvements of up to a factor of 2. As of 2020, the project is named Gallium.Nine. It is available now as a separate standalone package and no longer needs a patched Wine version. User interface Wine is usually invoked from the command-line interpreter: wine program.exe. winecfg There is the utility winecfg that starts a graphical user interface with controls for adjusting basic options. It is a GUI configuration utility included with Wine. Winecfg makes configuring Wine easier by making it unnecessary to edit the registry directly, although, if needed, this can be done with the included registry editor (similar to Windows regedit). Third-party applications Some applications require more tweaking than simply installing the application in order to work properly, such as manually configuring Wine to use certain Windows DLLs. The Wine project does not integrate such workarounds into the Wine codebase, instead preferring to focus solely on improving Wine's implementation of the Windows API. While this approach focuses Wine development on long-term compatibility, it makes it difficult for users to run applications that require workarounds. Consequently, many third-party applications have been created to ease the use of those applications that do not work out of the box within Wine itself. The Wine wiki maintains a page of current and obsolete third-party applications. Winetricks is a script to install some basic components (typically Microsoft DLLs and fonts) and tweak settings required for some applications to run correctly under Wine. It can fully automate the install of a number of applications and games, including applying any needed workarounds. Winetricks has a GUI. The Wine project will accept bug reports for users of Winetricks, unlike most third-party applications. It is maintained by Wine developer Austin English. Q4Wine is an open GUI for advanced setup of Wine. Wine-Doors is an application management tool for the GNOME desktop which adds functionality to Wine. Wine-Doors is an alternative to WineTools which aims to improve upon WineTools' features and extend on the original idea with a more modern design approach. IEs4Linux is a utility to install all versions of Internet Explorer, including versions 4 to 6 and version 7 (in beta). Wineskin is a utility to manage Wine engine versions and create wrappers for macOS. PlayOnLinux is an application to ease the installation of Windows applications (primarily games). There is also a corresponding Macintosh version called PlayOnMac. Lutris is an open-source application to install Windows games on Linux. Bordeaux is a proprietary Wine GUI configuration manager that runs winelib applications. It also supports installation of third-party utilities, installation of applications and games, and the ability to use custom configurations. Bordeaux currently runs on Linux, FreeBSD, PC-BSD, Solaris, OpenSolaris, OpenIndiana, and macOS computers. Bottles is an open-source graphical Wine prefix and runners manager for Wine based on GTK4+Libadwaita. It provides a repository-based dependency installation system and bottle versioning to restore a previous state. WineGUI is a free and open-source graphical interface to manage Wine. It allows a user to create Wine bottles and install Windows applications or games. Functionality The developers of the Direct3D portions of Wine have continued to implement new features such as pixel shaders to increase game support. Wine can also use native DLLs directly, thus increasing functionality, but then a license for Windows is needed unless the DLLs were distributed with the application itself. Wine also includes its own open-source implementations of several Windows programs, such as Notepad, WordPad, Control Panel, Internet Explorer, and Windows Explorer. The Wine Application Database (AppDB) is a community-maintained on-line database about which Windows programs works with Wine and how well they work. Backward compatibility Wine ensures good backward compatibility with legacy Windows applications, including those written for Windows 3.1x. Wine can mimic different Windows versions required for some programs, going as far back as Windows 2.0. However, Windows 1.x and Windows 2.x support was removed from Wine development version 1.3.12. If DOSBox is installed on the system (see below on MS-DOS), Wine development version 1.3.12 and later nevertheless show the "Windows 2.0" option for the Windows version to mimic, but Wine still will not run most Windows 2.0 programs because MS-DOS and Windows functions are not currently integrated. Backward compatibility in Wine is generally superior to that of Windows, as newer versions of Windows can force users to upgrade legacy Windows applications, and may break unsupported software forever as there is nobody adjusting the program for the changes in the operating system. In many cases, Wine can offer better legacy support than newer versions of Windows with "Compatibility Mode". Wine can run 16-bit Windows programs (Win16) on a 64-bit operating system, which uses an x86-64 (64-bit) CPU, a functionality not found in 64-bit versions of Microsoft Windows. WineVDM allows 16-bit Windows applications to run on 64-bit versions of Windows. Wine partially supports Windows console applications, and the user can choose which backend to use to manage the console (choices include raw streams, curses, and user32). When using the raw streams or curses backends, Windows applications will run in a Unix terminal. 64-bit applications Preliminary support for 64-bit Windows applications was added to Wine 1.1.10, in December 2008. , the support is considered stable. The two versions of Wine are built separately, and as a result only building wine64 produces an environment only capable of running x86-64 applications. , Wine has stable support for a WoW64 build, which allows both 32-bit and 64-bit Windows applications to run inside the same Wine instance. To perform such a build, one must first build the 64-bit version, and then build the 32-bit version referencing the 64-bit version. Just like Microsoft's WoW64, the 32-bit build process will add parts necessary for handling 32-bit programs to the 64-bit build. This functionality is seen from at least 2010. MS-DOS Early versions of Microsoft Windows run on top of MS-DOS, and Windows programs may depend on MS-DOS programs to be usable. Wine does not have good support for MS-DOS, but starting with development version 1.3.12, Wine tries running MS-DOS programs in DOSBox if DOSBox is available on the system. However, due to a bug, current versions of Wine incorrectly identify Windows 1.x and Windows 2.x programs as MS-DOS programs, attempting to run them in DOSBox (which does not work). Winelib Wine provides Winelib, which allows its shared-object implementations of the Windows API to be used as actual libraries for a Unix program. This allows for Windows code to be built into native Unix executables. Since October 2010, Winelib also works on the ARM platform. Non-x86 architectures Support for Solaris SPARC was dropped in version 1.5.26. ARM, Windows CE, and Windows RT Wine provides some support for ARM (as well as ARM64/AArch64) processors and the Windows flavors that run on it. , Wine can run ARM/Win32 applications intended for unlocked Windows RT devices (but not Windows RT programs). Windows CE support (either x86 or ARM) is missing, but an unofficial, pre-alpha proof-of-concept version called WineCE allows for some support. Wine for Android On 3 February 2013 at the FOSDEM talk in Brussels, Alexandre Julliard demonstrated an early demo of Wine running on Google's Android operating system. Experimental builds of WINE for Android (x86 and ARM) were released in late 2017. It has been routinely updated by the official developers ever since. The default builds do not implement cross-architecture emulation via QEMU, and as a result ARM versions will only run ARM applications that use the Win32 API. Microsoft applications Wine, by default, uses specialized Windows builds of Gecko and Mono to substitute for Microsoft's Internet Explorer and .NET Framework. Wine has built-in implementations of JScript and VBScript. It is possible to download and run Microsoft's installers for those programs through winetricks or manually. Wine is not known to have good support for most versions of Internet Explorer (IE). Of all the reasonably recent versions, Internet Explorer 8 for Windows XP is the only version that reports a usable rating on Wine's AppDB, out-of-the-box. However Google Chrome gets a gold rating (as of Wine 5.5-staging), and Microsoft's IE replacement web browser Edge, is known to be based on that browser (after switching from Microsoft's own rendering engine). Winetricks offer auto-installation for Internet Explorer 6 through 8, so these versions can be reasonably expected to work with its built-in workarounds. An alternative for installing Internet Explorer directly is to use the now-defunct IEs4Linux. It is not compatible with the latest versions of Wine, and the development of IEs4Linux is inactive. Other versions of Wine The core Wine development aims at a correct implementation of the Windows API as a whole and has sometimes lagged in some areas of compatibility with certain applications. Direct3D, for example, remained unimplemented until 1998, although newer releases have had an increasingly complete implementation. CrossOver CodeWeavers markets CrossOver specifically for running Microsoft Office and other major Windows applications, including some games. CodeWeavers employs Alexandre Julliard to work on Wine and contributes most of its code to the Wine project under the LGPL. CodeWeavers also released a new version called CrossOver Mac for Intel-based Apple Macintosh computers on 10 January 2007. Unlike upstream wine, CrossOver is notably able to run on the x64-only versions of macOS, using a technique known as "wine32on64". As of 2012, CrossOver includes the functionality of both the CrossOver Games and CrossOver Pro lines therefore CrossOver Games and CrossOver Pro are no longer available as single products. CrossOver Games was optimized for running Windows video games. Unlike CrossOver, it didn't focus on providing the most stable version of Wine. Instead, experimental features are provided to support newer games. Proton On 21 August 2018, Valve announced a new variation of Wine, named Proton, designed to integrate with the Linux version of the company's Steam software (including Steam installations built into their Linux-based SteamOS operating system and Steam Machine computers). Valve's goal for Proton is to enable Steam users on Linux to play games which lack a native Linux port (particularly back-catalog games), and ultimately, through integration with Steam as well as improvements to game support relative to mainline Wine, to give users "the same simple plug-and-play experience" that they would get if they were playing the game natively on Linux. Proton entered public beta immediately upon being announced. Valve had already been collaborating with CodeWeavers since 2016 to develop improvements to Wine's gaming performance, some of which have been merged to the upstream Wine project. Some of the specific improvements incorporated into Proton include Vulkan-based Direct3D 9, 10, 11, and 12 implementations via vkd3d, DXVK, and D9VK multi-threaded performance improvements via esync, improved handling of fullscreen games, and better automatic game controller hardware support. Proton is fully open-source and available via GitHub. WINE@Etersoft The Russian company Etersoft has been developing a proprietary version of Wine since 2006. WINE@Etersoft supports popular Russian applications (for example, 1C:Enterprise by 1C Company). Other projects using Wine source code Other projects using Wine source code include: OTVDM, a 16-bit app compatibility layer for 64-bit Windows. ReactOS, a project to write an operating system compatible with Windows NT versions 5.x and up (which includes Windows 2000 and its successors) down to the device driver level. ReactOS uses Wine source code considerably; however due to architectural differences with ReactOS its code is not generally reused in Wine, such as in the case of ReactOS specific DLLs, such as ntdll, user32, kernel32, gdi32, and advapi. In July 2009, Aleksey Bragin, the ReactOS project lead, started a new ReactOS branch called Arwinss, and it was officially announced in January 2010. Arwinss is an alternative implementation of the core Win32 components, and uses mostly unchanged versions of Wine's user32.dll and gdi32.dll. WineBottler, a wrapper around Wine in the form of a normal Mac application. It manages multiple Wine configurations for different programs in the form of "bottles." Wineskin, an open source Wine GUI configuration manager for macOS. Wineskin creates a wrapper around Wine in the form of a normal Mac Application. The wrapper can also be used to make a distributable "port" of software. Odin, a project to run Win32 binaries on OS/2 or convert them to OS/2 native format. The project also provides the Odin32 API to compile Win32 programs for OS/2. Virtualization products such as Parallels Desktop for Mac and VirtualBox use WineD3D to make use of the GPU. WinOnX, a commercial package of Wine for macOS that includes a GUI for adding and managing applications and virtual machines. WineD3D for Windows, a compatibility wrapper which emulates old Direct3D versions and features that were removed by Microsoft in recent Windows releases, using OpenGL. This sometimes gets older games working again. Apple Game Porting Toolkit, a suite of software introduced at Apple's Worldwide Developer Conference in June 2023 to facilitate porting games from Windows to Mac. Discontinued Cedega / WineX: TransGaming Inc. (now Findev Inc. since the sale of its software businesses) produced the proprietary Cedega software. Formerly known as WineX, Cedega represented a fork from the last MIT-licensed version of Wine in 2002. Much like CrossOver Games, TransGaming's Cedega was targeted towards running Windows video games. On 7 January 2011, TransGaming Inc. announced continued development of Cedega Technology under the GameTree Developer Program. TransGaming Inc. allowed members to keep using their Cedega ID and password until 28 February 2011. Cider: TransGaming also produced Cider, a library for Apple–Intel architecture Macintoshes. Instead of being an end-user product, Cider (like Winelib) is a wrapper allowing developers to adapt their games to run natively on Intel Mac without any changes in source code. Darwine: a port of the Wine libraries to Darwin and Mac OS X for the PowerPC and Intel x86 (32-bit) architectures, created by the OpenDarwin team in 2004. Its PowerPC version relied on QEMU. Darwine was merged back into Wine in 2009. E/OS LX: a project attempting to allow any program designed for any operating system to be run without the need to actually install any other operating system. Pipelight: a custom version of Wine (wine-compholio) that acts as a wrapper for Windows NPAPI plugins within Linux browsers. This tool permits Linux users to run Microsoft Silverlight, the Microsoft equivalent of Adobe Flash, and the Unity web plugin, along with a variety of other NPAPI plugins. The project provides an extensive set of patches against the upstream Wine project, some of which were approved and added to upstream Wine. Pipelight is largely obsolete, as modern browsers no longer support NPAPI plugins and Silverlight has been deprecated by Microsoft. Reception The Wine project has received a number of technical and philosophical complaints and concerns over the years. Security Because of Wine's ability to run Windows binary code, concerns have been raised over native Windows viruses and malware affecting Unix-like operating systems as Wine can run limited malware made for Windows. A 2018 security analysis found that 5 out of 30 malware samples were able to successfully run through Wine, a relatively low rate that nevertheless posed a security risk. For this reason the developers of Wine recommend never running it as the superuser. Malware research software such as ZeroWine runs Wine on Linux in a virtual machine, to keep the malware completely isolated from the host system. An alternative to improve the security without the performance cost of using a virtual machine, is to run Wine in an LXC container, as Anbox software is doing by default with Android. Another security concern is when the implemented specifications are ill-designed and allow for security compromise. Because Wine implements these specifications, it will likely also implement any security vulnerabilities they contain. One instance of this problem was the 2006 Windows Metafile vulnerability, which saw Wine implementing the vulnerable SETABORTPROC escape. Wine vs. native Unix applications A common concern about Wine is that its existence means that vendors are less likely to write native Linux, macOS, and BSD applications. As an example of this, it is worth considering IBM's 1994 operating system, OS/2 Warp. An article describes the weaknesses of OS/2 which killed it, the first one being: However, OS/2 had many problems with end user acceptance. Perhaps the most serious was that most computers sold already came with DOS and Windows, and many people didn't bother to evaluate OS/2 on its merits due to already having an operating system. "Bundling" of DOS and Windows and the chilling effect this had on the operating system market frequently came up in United States v. Microsoft Corporation. The Wine project itself responds to the specific complaint of "encouraging" the continued development for the Windows API on one of its wiki pages: Also, the Wine Wiki page claims that Wine can help break the chicken-and-egg problem for Linux on the desktop: The use of Wine for gaming has proved specifically controversial in the Linux community, as some feel it is preventing, or at least hindering, the further growth of native Linux gaming on the platform. One quirk however is that Wine is now able to run 16-bit and even certain 32-bit applications and games that do not launch on current 64-bit Windows versions. This use-case has led to running Wine on Windows itself via Windows Subsystem for Linux or third-party virtual machines, as well as encapsulated by means such as BoxedWine and Otvdm. Microsoft Until 2020, Microsoft had not made any public statements about Wine. However, the Windows Update online service will block updates to Microsoft applications running in Wine. On 16 February 2005, Ivan Leo Puoti discovered that Microsoft had started checking the Windows Registry for the Wine configuration key and would block the Windows Update for any component. As Puoti noted: "It's also the first time Microsoft acknowledges the existence of Wine." In January 2020, Microsoft cited Wine as a positive consequence of being able to reimplement APIs, in its amicus curiae brief for Google LLC v. Oracle America, Inc. In August 2024, Microsoft donated the Mono Project, a reimplementation of the .NET Framework, to the developers of Wine.	Technology	System	null
37595	https://en.wikipedia.org/wiki/Weasel	Weasel	Weasels are mammals of the genus Mustela of the family Mustelidae. The genus Mustela includes the least weasels, polecats, stoats, ferrets, and European mink. Members of this genus are small, active predators, with long and slender bodies and short legs. The family Mustelidae, or mustelids (which also includes badgers, otters, and wolverines), is often referred to as the "weasel family". In the UK, the term "weasel" usually refers to the smallest species, the least weasel (M. nivalis), the smallest carnivoran species. Least weasels vary in length from , females being smaller than the males, and usually have red or brown upper coats and white bellies; some populations of some species moult to a wholly white coat in winter. They have long, slender bodies, which enable them to follow their prey into burrows. Their tails may be from long. Weasels feed on small mammals and have from time to time been considered vermin because some species took poultry from farms or rabbits from commercial warrens. They do, on the other hand, eat large numbers of rodents. Their range spans Europe, North America, much of Asia and South America, and small areas in North Africa. Terminology The English word "weasel" was originally applied to one species of the genus, the European form of the least weasel (Mustela nivalis). This usage is retained in British English, where the name is also extended to cover several other small species of the genus. However, in technical discourse and in American usage, the term "weasel" can refer to any member of the genus, the genus as a whole, and even to members of the related genus Neogale. Of the 16 extant species currently classified in the genus Mustela, 10 have "weasel" in their common names. Among those that do not are the three species of ermine, the polecats, the ferret, and the European mink. The American mink and the extinct sea mink were commonly included in this genus as Mustela vison and Mustela macrodon, respectively, but in 1999 they were moved to the genus Neovison. In 2021, both Neovison species, along with the long-tailed weasel (Mustela frenata), Amazon weasel (Mustela africana) and Colombian weasel (Mustela felipei) were moved to the genus Neogale, as the clade containing these five species was found to be fully distinct from Mustela. Taxonomy The genus name Mustela comes from the Latin word for weasel combining the words mus meaning "mouse" and telum meaning "javelin" for its long body. Species The following information is according to the Integrated Taxonomic Information System and MammalDiversity. 1 Europe and Northern Asia division excludes China. Cultural meanings Weasels have been assigned a variety of cultural meanings. In Greek culture, a weasel near one's house is a sign of bad luck, even evil, "especially if there is in the household a girl about to be married", since the animal (based on its Greek etymology) was thought to be an unhappy bride who was transformed into a weasel and consequently delights in destroying wedding dresses. In Macedonia, however, weasels are generally seen as an omen of good fortune. In early-modern Mecklenburg, Germany, amulets from weasels were deemed to have strong magic; the period between 15 August and 8 September was specifically designated for the killing of weasels. In Montagne Noire (France), Ruthenia, and the early medieval culture of the Wends, weasels were not meant to be killed. According to Daniel Defoe also, meeting a weasel is a bad omen. In English-speaking areas, weasel can be an insult, noun or verb, for someone regarded as sneaky, conniving or untrustworthy. Similarly, "weasel words" is a critical term for words or phrasing that are vague, misleading or equivocal. Japanese superstitions In Japan, were seen as yōkai (causing strange occurrences). According to the encyclopedia Wakan Sansai Zue from the Edo period, a pack of weasels would cause conflagrations, and the cry of a weasel was considered a harbinger of misfortune. In the Niigata Prefecture, the sound of a pack of weasels making a rustle resembled six people hulling rice, so was called the "weasel's six-person mortar", and it was an omen for one's home to decline or flourish. It is said that when people chase after this sound, the sound stops. They are also said to shapeshift like the fox (kitsune) or tanuki, and the nyūdō-bōzu told about in legends in the Tōhoku region and the Chūbu region are considered weasels in disguise, and they are also said to shapeshift into ōnyūdō and little monks. In the collection of depictions Gazu Hyakki Yagyō by Sekien Toriyama, they were depicted under the title 鼬, but they were read not as "itachi", but rather as "ten", and "ten" were considered to be weasels that have reached one hundred years of age and became yōkai that possessed supernatural powers. Another theory is that when weasels reach several hundred years of age, they become mujina (Japanese badgers). In Japanese, weasels are called and in the Tōhoku Region and Shinshu, it was believed that there were families that were able to use a certain practice to freely use kudagitsune as iizuna-tsukai or kitsune-mochi. It is said that Mount Iizuna, from the Nagano Prefecture, got its name due to how the gods gave people mastery of this technique from there. According to the folklorist Mutō Tetsujō, "They are called izuna in the Senboku District, Akita Prefecture, and there are also the ichiko (itako) that use them." Also, in the Kitaakita District, they are called mōsuke (猛助), and they are feared as yōkai even more than foxes (kitsune). In the Ainu language, ermines are called upas-čironnup or sáčiri, but since least weasels are also called sáčiri, Mashio Chiri surmised that the honorary title poy-sáčiri-kamuy (where poy means "small") refers to least weasels. Kamaitachi Kamaitachi is a phenomenon wherein one who is idle is suddenly injured as if his or her skin were cut by a scythe. In the past, this was thought to be "the deed of an invisible yōkai weasel". An alternate theory, asserts that kamaitachi is derived from , so were not originally related to weasels at all.	Biology and health sciences	Carnivora	null
37602	https://en.wikipedia.org/wiki/Eagle	Eagle	Eagle is the common name for the golden eagle, bald eagle, and other birds of prey in the family of the Accipitridae. Eagles belong to several groups of genera, some of which are closely related. True eagles comprise the genus Aquila. Most of the 68 species of eagles are from Eurasia and Africa. Outside this area, just 14 species can be found—two in North America, nine in Central and South America, and three in Australia. Eagles are not a natural group but denote essentially any kind of bird of prey large enough to hunt sizeable (about 50 cm long or more overall) vertebrates. Etymology The word "eagle" is borrowed into English from and , both derived ultimately from ("eagle"). It is cognate with terms such as , and . It is broadly synonymous with the less common English term "erne" or "earn", deriving from , from , in which it acts as the usual word for the bird. The Old English term is turn derived from and is cognate with other synonymous words in Germanic languages such as , and . Through the Proto-Indo-European root, it is further related to words such as ("bird") and ("eagle"). Although "erne" can be used to refer to any eagle, it is most commonly used for the golden eagle or sea-eagle. Description Eagles are large, powerfully-built birds of prey, with heavy heads and beaks. Even the smallest eagles, such as the booted eagle (Hieraaetus pennatus), which is comparable in size to a common buzzard (Buteo buteo) or red-tailed hawk (B. jamaicensis), have relatively longer and more evenly broad wings, and more direct, faster flight, despite the reduced size of their aerodynamic feathers. Most eagles are larger than any other raptors, apart from some vultures. The smallest species of eagle is the Great Nicobar serpent eagle (Spilornis klossi), at and . The largest species are discussed below. Like all birds of prey, eagles have very large hooked beaks for ripping flesh from their prey, strong, muscular legs, and powerful talons. The beak is typically heavier than that of most other birds of prey. Eagles' eyes are extremely powerful. It is estimated that the wedge-tailed eagle has a visual acuity twice that of a typical human. This acuity enables eagles to spot potential prey from a very long distance. This keen eyesight is primarily attributed to their extremely large pupils which ensure minimal diffraction (scattering) of the incoming light. Like most diurnal raptors, eagles have little ability to see ultraviolet light. The female of all known species of eagles is larger than the male. Eagles normally build their nests, called eyries, in tall trees or on high cliffs. Many species lay two eggs, but the older, larger chick frequently kills its younger sibling once it has hatched. The parents take no action to stop the killing. It is said that eagles fly above clouds but this is not true. Eagles fly during storms and glide from the wind's pressure. This saves the bird's energy. Due to the size and power of many eagle species, they are ranked at the top of the food chain as apex predators in the avian world. The type of prey varies by genus. The Haliaeetus and Icthyophaga eagles prefer to capture fish, though the species in the former often capture various animals, especially other water birds, and are powerful kleptoparasites of other birds. The snake and serpent eagles of the genera Circaetus, Terathopius, and Spilornis predominantly prey on the great diversity of snakes found in the tropics of Africa and Asia. The eagles of the genus Aquila are often the top birds of prey in open habitats, taking almost any medium-sized vertebrate they can catch. Where Aquila eagles are absent, other eagles, such as the buteonine black-chested buzzard-eagle of South America, may assume the position of top raptorial predator in open areas. Many other eagles, including the species-rich genus Spizaetus, live predominantly in woodlands and forests. These eagles often target various arboreal or ground-dwelling mammals and birds, which are often unsuspectingly ambushed in such dense, knotty environments. Hunting techniques differ among the species and genera, with some individual eagles having engaged in quite varied techniques based on their environment and prey at any given time. Most eagles grab prey without landing and take flight with it, so the prey can be carried to a perch and torn apart. The bald eagle is noted for having flown with the heaviest load verified to be carried by any flying bird, since one eagle flew with a mule deer fawn. However, a few eagles may target prey considerably heavier than themselves; such prey is too heavy to fly with, thus it is either eaten at the site of the kill or taken in pieces back to a perch or nest. Golden and crowned eagles have killed ungulates weighing up to and a martial eagle even killed a duiker, 7–8 times heavier than the preying eagle. Authors on birds David Allen Sibley, Pete Dunne, and Clay Sutton described the behavioral difference between hunting eagles and other birds of prey thus (in this case the bald and golden eagles as compared to other North American raptors): They have at least one singular characteristic. It has been observed that most birds of prey look back over their shoulders before striking prey (or shortly thereafter); predation is after all a two-edged sword. All hawks seem to have this habit, from the smallest kestrel to the largest Ferruginous – but not the Eagles. Among the eagles are some of the largest birds of prey: only the condors and some of the Old World vultures are markedly larger. It is regularly debated which should be considered the largest species of eagle. They could be measured variously in total length, body mass, or wingspan. Different lifestyle needs among various eagles result in variable measurements from species to species. For example, many forest-dwelling eagles, including the very large harpy eagle, have relatively short wingspans, a feature necessary for being able to maneuver in quick, short bursts through densely forested habitats. Eagles in the genus Aquila, found almost exclusively in open country, are noted for their ability to soar, and have relatively long wings for their size. These lists of the top five eagles are based on weight, length, and wingspan, respectively. Unless otherwise noted by reference, the figures listed are the median reported for each measurement in the guide Raptors of the World in which only measurements that could be personally verified by the authors were listed. Habitat The eagles are generally distributed in all types of habitats and nearly all parts of the world. The birds can be found in northern tundra to tropical rainforests and deserts. In North America, bald eagles and golden eagles are very common. Distribution Australasian Australia: wedge-tailed eagle (range extends into southern New Guinea), white-bellied sea-eagle (range extends into Asia), little eagle. New Guinea: Papuan eagle, white-bellied sea-eagle, pygmy eagle. Nearctic (USA and Canada): golden eagle (also found in Palearctic), bald eagle. Neotropical (Central and South America): Spizaetus (four species), solitary eagles (two spp.), harpy eagle, crested eagle, black-chested buzzard-eagle. Palearctic (Europe, Northern Africa, Asia without South Asia and Southeast Asia) Eurasia: Golden eagle, White-tailed eagle. Subsaharan Africa: African fish eagle, Martial Eagle, Crowned eagle, Verreaux's eagle, Tawny eagle, Long-crested eagle Groups Eagles are often informally divided into four groups. The snake eagles are placed in the subfamily Circaetinae. The fish eagles, booted eagles, and harpy eagles have traditionally been placed in the subfamily Buteoninae together with the buzzard-hawks (buteonine hawks) and harriers. Some authors may treat these groups as tribes of the Buteoninae; Lerner & Mindell proposed separating the eagle groups into their own subfamilies of Accipitridae. Fish eagles Sea eagles or fish eagles take fish as a large part of their diets, either fresh or as carrion. Proposed subfamily Haliaeetinae. Genera: Haliaeetus, Icthyophaga. Some authors include Gypohierax angolensis, the "vulturine fish eagle" (also called the palm-nut vulture) in this group. However, genetic analyses indicate it is related to a grouping of Neophron–Gypaetus–Eutriorchis (Egyptian vulture, bearded vulture (lammergeier), and Madagascar serpent eagle). The fish eagles have a close genetic relationship with Haliastur and Milvus; the whole group is only distantly related to the Buteo group. Fish eagles exist in every continent throughout the world, except for South America. Although fish eagles can be found in many different places around the world, they have been classified as "Near Threatened". Reasons such as overfishing, pollution, habitat destruction, and the use of pesticides have contributed to the species' rapid population drop. Booted eagles The booted eagle is a group of eagle that typically migrates across the Sahara Desert to Europe. It usually reaches Europe around the beginning of March and leaves by the end of September. It's interesting to note that these types of eagles usually mate with the same partner and return to the same areas years later. Female booted eagles usually lay 1-4 eggs, which promptly hatch after 37 to 40 days. Researchers estimate that there are between 3600 and 6900 pairs of booted eagles in Europe, which are mostly situated in the Iberian Peninsula. Booted eagles or "true eagles" have feathered tarsi (lower legs). Tribe Aquililae or proposed subfamily Aquilinae. Genera: Aquila, Hieraaetus; Spizaetus, Oroaetus, Spizastur; Nisaetus; Ictinaetus, Lophoaetus; Polemaetus; and Stephanoaetus. See comments under eagle species for changes to the composition of these genera. Snake eagles Most snake or serpent eagles, as the name suggests, primarily prey on snakes. Subfamily Circaetinae. Genera: Circaetus, Spilornis, Dryotriorchis, Terathopius. Eutriorchis (subfamily Gypaetinae or Circaetinae). Despite filling the niche of a snake eagle, genetic studies suggest that the Madagascar serpent eagle (Eutriorchis) is not related to them. Over several decades, a great deal of research has been done on the Snake-eagle's diet, which is mainly made up of reptiles, especially snakes. When it comes to catching snakes, it is generally accepted that the bird exhibits generalist feeding behavior, which means it does not hunt down specific types of snakes but rather feeds on them depending on their availability in the wild. Harpy eagles Harpy eagles or "giant forest eagles" are large eagles that inhabit tropical forests. The group contains two to six species, depending on the author. Although these birds occupy similar niches and have traditionally been grouped, they are not all related: the solitary eagles are related to the black hawks and the Philippine eagle to the snake eagles. Harpy eagles (proposed subfamily Harpiinae) Harpia harpyja, harpy eagle ― Central and South America. Morphnus guianensis, crested eagle ― Central and South America. Harpyopsis novaeguineae, Papuan eagle ― New Guinea. Philippine eagle Pithecophaga jefferyi, Philippine eagle ― Philippines. Solitary eagles Chaco eagle or crowned solitary eagle, Buteogallus (formerly Harpyhaliaetus) coronatus ― South America. Solitary eagle or montane solitary eagle, Buteogallus (formerly Harpyhaliaetus) solitarius ― South America. Species Major new research into eagle taxonomy suggests that the important genera Aquila and Hieraaetus are not composed of nearest relatives, and it is likely that a reclassification of these genera will soon take place, with some species being moved to Lophaetus or Ictinaetus. Bonelli's eagle and the African hawk-eagle have been moved from Hieraaetus to Aquila. Either the greater spotted eagle and lesser spotted eagle should move from Aquila to join the long-crested eagle in Lophaetus, or, perhaps better, all three of these species should move to Ictinaetus with the black eagle. The steppe eagle and tawny eagle, once thought to be conspecific, are not even each other's nearest relatives. Family Accipitridae Subfamily Buteoninae – hawks (buzzards), true eagles and seaeagles Genus Geranoaetus Black-chested buzzard-eagle, Geranoaetus melanoleucus Genus Harpyhaliaetus Chaco eagle, Buteogallus coronatus Solitary eagle, H. solitarius Genus Morphnus Crested eagle, Morphnus guianensis Genus Harpia Harpy eagle, Harpia harpyja Genus Pithecophaga Philippine eagle, Pithecophaga jefferyi Genus Harpyopsis Papuan eagle, Harpyopsis novaeguineae Genus Spizaetus Black hawk-eagle, S. tyrannus Ornate hawk-eagle, S. ornatus Black-and-white hawk-eagle, S. melanoleucus – formerly Spizastur Black-and-chestnut eagle, S. isidori – formerly Oroaetus Genus Nisaetus – previously included in Spizaetus Changeable hawk-eagle, N. cirrhatus Flores hawk-eagle N. floris – earlier a subspecies, S. c. floris Sulawesi hawk-eagle, N. lanceolatus Mountain hawk-eagle, N. nipalensis Legge's hawk-eagle, Nisaetus kelaarti – previously a race of S. nipalensis Blyth's hawk-eagle, N. alboniger Javan hawk-eagle, N. bartelsi (Northern) Philippine hawk-eagle, N. philippensis Pinsker's hawk-eagle (Southern Philippine hawk-eagle), Nisaetus pinskeri – earlier S. philippensis pinskeri Wallace's hawk-eagle, N. nanus Genus Lophaetus Long-crested eagle, Lophaetus occipitalis – possibly belongs in Ictinaetus Genus Stephanoaetus Crowned eagle, Stephanoaetus coronatus Malagasy crowned eagle, Stephanoaetus mahery Genus Polemaetus Martial eagle, Polemaetus bellicosus Genus Hieraaetus Ayres's hawk-eagle, H. ayresii Little eagle, H. morphnoides Pygmy eagle, H. weiskei – previously subspecies H. m. weiskei Booted eagle, H. pennatus Haast's eagle, †H. moorei Genus Lophotriorchis Rufous-bellied eagle, L. kienerii Genus Aquila Bonelli's eagle, Aquila fasciata – formerly Hieraaetus fasciatus African hawk-eagle, A. spilogaster – formerly in Hieraaetus Cassin's hawk-eagle, A. africana – formerly in Hieraaetus or Spizaetus genera Golden eagle, A. chrysaetos Eastern imperial eagle, A. heliaca Spanish imperial eagle A. adalberti Steppe eagle, A. nipalensis Tawny eagle, A. rapax Greater spotted eagle, A. clanga – to be moved to Lophaetus or Ictinaetus Lesser spotted eagle, A. pomarina – to be moved to Lophaetus or Ictinaetus Indian spotted eagle, A. hastata – to be moved to Lophaetus or Ictinaetus Verreaux's eagle, A. verreauxii Gurney's eagle, A. gurneyi Wahlberg's eagle, A. wahlbergi – to be moved to Hieraaetus Wedge-tailed eagle, A. audax Genus Ictinaetus Black eagle, Ictinaetus malaiensis Genus Haliaeetus White-tailed eagle, Haliaeetus albicilla Bald eagle, H. leucocephalus Steller's sea eagle, H. pelagicus Pallas' sea eagle, H. leucoryphus Genus Icthyophaga Lesser fish eagle, Icthyophaga humilis Grey-headed fish eagle, I. ichthyaetus African fish eagle, I. vocifer White-bellied sea eagle, I. leucogaster Sanford's sea eagle, I. sanfordi Madagascar fish eagle, I. vociferoides Subfamily Circaetinae: snake-eagles Genus Terathopius Bateleur, Terathopius ecaudatus Genus Circaetus Short-toed snake eagle, Circaetus gallicus Beaudouin's snake eagle, Circaetus beaudouini Black-chested snake eagle, C. pectoralis Brown snake eagle, C. cinereus Fasciated snake eagle, C. fasciolatus Western banded snake eagle, C. cinerascens Genus Dryotriorchis Congo serpent eagle, D. spectabilis Genus Spilornis Crested serpent eagle, Spilornis cheela Central Nicobar serpent eagle, S. minimus (subspecies or species) Great Nicobar serpent eagle, S. klossi Mountain serpent eagle, S. kinabaluensis Sulawesi serpent eagle, S. rufipectus Philippine serpent eagle, S. holospilus Andaman serpent eagle, S. elgini Genus Eutriorchis Madagascar serpent eagle, Eutriorchis astur In culture Etymology The modern English term for the bird is derived from by way of . The origin of is unknown, but it is believed to possibly derive from (meaning dark-colored, swarthy, or blackish) as a reference to the plumage of eagles. Old English used the term , related to Scandinavia's ørn/örn. It is similar to other Indo-European terms for "bird" or "eagle", including (), (), and . In the southern part of Finland, near the Gulf of Finland, is the town of Kotka, which literally means "eagle", while the town of L'Aquila in the central part of Italy literally means "the eagle". In Britain before 1678, eagle referred specifically to the golden eagle, with the other native species, the white-tailed eagle, being known as erne. The modern name "golden eagle" for aquila chrysaetos was introduced by the naturalist John Ray. The village of Eagle in Lincolnshire, England, has nothing to do with the bird; its name is derived from the Old English words for "oak" and "wood" (compare Oakley). Religion and spirituality In the ancient Sumerian mythology, the mythical king Etana was said to have been carried into heaven by an eagle. Classical writers such as Lucan and Pliny the Elder claimed that the eagle was able to look directly at the sun, and that they forced their fledglings to do the same. Those that blinked would be cast from the nest. This belief persisted until the Medieval era. The eagle is the patron animal of the ancient Greek god Zeus. In particular, Zeus was said to have taken the form of an eagle in order to abduct Ganymede, and there are numerous artistic depictions of the eagle Zeus bearing Ganymede aloft, from Classical times up to the present (see illustrations in the Ganymede (mythology) page.) Eagles appear metaphorically in many translations of the Old Testament. God is spoken of as carrying Israel on "eagles' wings" in Exodus 19:4, Isaiah 40:31 compares those who wait on the Lord to flying eagles, and Psalm 103 mentions renewing one's youth "as the eagle". In explaining this rejuvenation, Augustine of Hippo says in his commentary on the Psalms that eagles' beaks overgrow as they age and that they break them against rocks to restore them. The translation, however, is uncertain: the word in the Hebrew, נשר, can also be translated vulture, and is listed alongside specific kinds of vulture in Leviticus' discussion of unclean animals. The eagle is also often used in Christian iconography to represent the Gospel of John, and eagle-shaped lecterns are common in Anglican and some Roman Catholic churches. The eagle was believed to be able to look directly into the sun in the same way that the Gospel of John looks directly at Jesus' divinity, and the great distances the eagle flies represent the spread of the gospel to the ends of the earth. The United States eagle feather law stipulates that only individuals of certifiable Native American ancestry enrolled in a federally recognized tribe are legally authorized to obtain eagle feathers for religious or spiritual reasons. In Canada, the poaching of eagle feathers for the booming U.S. market has sometimes resulted in the arrests of First Nations person for the crime. The Moche people of ancient Peru worshiped the eagle and often depicted eagles in their art. The golden eagle was sacred to the Aztec god Huitzilopochtli while the harpy eagle was sacred to Quetzalcoatl. Heraldry Eagles are an exceptionally common symbol in heraldry, being considered the "King of Birds" in contrast to the lion, the "King of Beasts". Whereas the lion (e.g. England) usually represents authority, the eagle is the symbol of power. They are particularly popular in Germanic countries such as Austria, due to their association with the Holy Roman Empire. The eagle of the Holy Roman Empire was two-headed, supposedly representing the two divisions, East and West, of the old Roman Empire. This motif, derived from the Byzantine (Eastern Roman) Empire was also adopted by the Russian Empire and is still featured in the Flag of Albania. The Roman eagle was preceded by the eagle of Ptolemaic Egypt and the Achaemenid Empire. In the coat of arms of Kotka, Finland, the eagle is depicted carrying an anchor and the caduceus on its feet. Heraldic eagles are most often found displayed, i.e. with their wings and legs extended. They can also occur close, i.e. with their wings folded, or rising, i.e. about to take flight. The heads, wings, and legs of eagles can also be found independently. Eagles symbolize strength, courage, and independence and are commonly found in the heraldry of many nations across the world. Albania, Andorra, Armenia, Austria, Dagestan, Egypt, Germany, Ghana, Indonesia, Iraq, Jordan, Kazakhstan, Mexico, Montenegro, Nigeria, Philippines, Poland, Palestine, Panama, Russia, Romania, Serbia, South Sudan, Somaliland, the United States of America, Yemen, Zambia, and Zimbabwe are the nations whose coats of arms feature an eagle. The eagle's continuing significance and worldwide appeal as a forceful symbol in national identity and imagery is demonstrated by its widespread usage.	Biology and health sciences	Accipitriformes and Falconiformes	null
37618	https://en.wikipedia.org/wiki/Finger	Finger	A finger is a prominent digit on the forelimbs of most tetrapod vertebrate animals, especially those with prehensile extremities (i.e. hands) such as humans and other primates. Most tetrapods have five digits (pentadactyly), and short digits (i.e. significantly shorter than the metacarpal/metatarsals) are typically referred to as toes, while those that are notably elongated are called fingers. In humans, the fingers are flexibly articulated and opposable, serving as an important organ of tactile sensation and fine movements, which are crucial to the dexterity of the hands and the ability to grasp and manipulate objects. Land vertebrate fingers As terrestrial vertebrates were evolved from lobe-finned fish, their forelimbs are phylogenetically equivalent to the pectoral fins of fish. Within the taxa of the terrestrial vertebrates, the basic pentadactyl plan, and thus also the metacarpals and phalanges, undergo many variations. Morphologically the different fingers of terrestrial vertebrates are homolog. The wings of birds and those of bats are not homologous, they are analogue flight organs. However, the phalanges within them are homologous. Chimpanzees have lower limbs that are specialized for manipulation, and (arguably) have fingers (instead of toes) on their lower limbs as well. In the case of primates in general, the digits of the hand are overwhelmingly referred to as "fingers". Primate fingers have both fingernails and fingerprints. Research has been carried out on the embryonic development of domestic chickens showing that an interdigital webbing forms between the tissues that become the toes, which subsequently regresses by apoptosis. If apoptosis fails to occur, the interdigital skin remains intact. Many animals have developed webbed feet or skin between the fingers from this like the Wallace's flying frog. Human fingers Usually humans have five digits, the bones of which are termed phalanges, on each hand, although some people have more or fewer than five due to congenital disorders such as polydactyly or oligodactyly, or accidental or intentional amputations. The first digit is the thumb, followed by the index finger, middle finger, ring finger, and little finger or pinkie. According to different definitions, the thumb can be called a finger, or not. English dictionaries describe finger as meaning either one of the five digits including the thumb, or one of the four digits excluding the thumb (in which case they are numbered from 1 to 4 starting with the index finger closest to the thumb). Structure Skeleton The thumb (connected to the trapezium) is located on one of the sides, parallel to the arm. The palm has five bones known as metacarpal bones, one to each of the five digits. Human hands contain fourteen digital bones, also called phalanges, or phalanx bones: two in the thumb (the thumb has no middle phalanx) and three in each of the four fingers. These are the distal phalanx, carrying the nail, the middle phalanx, and the proximal phalanx. Joints are formed wherever two or more of these bones meet. Each of the fingers has three joints: metacarpophalangeal joint (MCP) – the joint at the base of the finger proximal interphalangeal joint (PIP) – the joint in the middle of the finger distal interphalangeal joint (DIP) – the joint closest to the fingertip. Sesamoid bones are small ossified nodes embedded in the tendons to provide extra leverage and reduce pressure on the underlying tissue. Many exist around the palm at the bases of the digits; the exact number varies between different people. The articulations are: interphalangeal articulations between phalangeal bones, and metacarpophalangeal joints connecting the phalanges to the metacarpal bones. Muscles Each finger may flex and extend, abduct and adduct, and so also circumduct. Flexion is by far the strongest movement. In humans, there are two large muscles that produce flexion of each finger, and additional muscles that augment the movement. The muscle bulks that move each finger may be partly blended, and the tendons may be attached to each other by a net of fibrous tissue, preventing completely free movement. Although each finger seems to move independently, moving one finger also moves the other fingers slightly which is called finger interdependence or finger enslaving. Fingers do not contain muscles (other than arrector pili). The muscles that move the finger joints are in the palm and forearm. The long tendons that deliver motion from the forearm muscles may be observed to move under the skin at the wrist and on the back of the hand. Muscles of the fingers can be subdivided into extrinsic and intrinsic muscles. The extrinsic muscles are the long flexors and extensors. They are called extrinsic because the muscle belly is located on the forearm. The fingers have two long flexors, located on the underside of the forearm. They insert by tendons to the phalanges of the fingers. The deep flexor attaches to the distal phalanx, and the superficial flexor attaches to the middle phalanx. The flexors allow for the actual bending of the fingers. The thumb has one long flexor and a short flexor in the thenar muscle group. The human thumb also has other muscles in the thenar group (opponens and abductor brevis muscle), moving the thumb in opposition, making grasping possible. The extensors are located on the back of the forearm and are connected in a more complex way than the flexors to the dorsum of the fingers. The tendons unite with the interosseous and lumbrical muscles to form the extensorhood mechanism. The primary function of the extensors is to straighten out the digits. The thumb has two extensors in the forearm; the tendons of these form the anatomical snuff box. Also, the index finger and the little finger have an extra extensor, used for instance for pointing. The extensors are situated within six separate compartments. The first compartment contains abductor pollicis longus and extensor pollicis brevis. The second compartment contains extensors carpi radialis longus and brevis. The third compartment contains extensor pollicis longus. The extensor digitorum indicis and extensor digitorum communis are within the fourth compartment. Extensor digiti minimi is in the fifth, and extensor carpi ulnaris is in the sixth. The intrinsic muscle groups are the thenar and hypothenar muscles (thenar referring to the thumb, hypothenar to the small finger), the dorsal and palmar interossei muscles (between the metacarpal bones) and the lumbrical muscles. The lumbricals arise from the deep flexor (and are special because they have no bony origin) and insert on the dorsal extensor hood mechanism. Skin Aside from the genitals, the fingertips possess the highest concentration of touch receptors and thermoreceptors among all areas of the human skin, making them extremely sensitive to temperature, pressure, vibration, texture and moisture. A study in 2013 suggested fingers can feel nano-scale wrinkles on a seemingly smooth surface, a level of sensitivity not previously recorded. This makes the fingers commonly used sensory probes to ascertain properties of objects encountered in the world, making them prone to injury. The of a finger is the fleshy mass on the palmar aspect of the extremity of the finger. Fingertip wrinkling in water Although a common phenomenon, the underlying functions and mechanism of fingertip wrinkling following immersion in water are relatively unexplored. Originally it was assumed that the wrinkles were simply the result of the skin swelling in water, but it is now understood that the furrows are caused by the blood vessels constricting due to signalling by the sympathetic nervous system in response to water exposure. One hypothesis for why this occurs, the "rain tread" hypothesis, posits that the wrinkles may help the fingers grip things when wet, possibly being an adaption from a time when humans dealt with rain and dew in forested primate habitats. A 2013 study supporting this hypothesis found that the wrinkled fingertips provided better handling of wet objects but gave no advantage for handling dry objects. However, a 2014 study attempting to reproduce these results was unable to demonstrate any improvement of handling wet objects with wrinkled fingertips. Regrowth of the fingertips Fingertips, after having been torn off children, have been observed to regrow in less than 8 weeks. However, these fingertips do not look the same, although they do look more appealing than a skin graft or a sewn fingertip. No healing occurs if the tear happens below the nail. This works because the distal phalanges are regenerative in youth, and stem cells in the nails create new tissue that ends up as the fingertip. Brain representation Each finger has an orderly somatotopic representation on the cerebral cortex in the somatosensory cortex area 3b, part of area 1 and a distributed, overlapping representations in the supplementary motor area and primary motor area. The somatosensory cortex representation of the hand is a dynamic reflection of the fingers on the external hand: in syndactyly people have a clubhand of webbed, shortened fingers. However, not only are the fingers of their hands fused, but the cortical maps of their individual fingers also form a club hand. The fingers can be surgically divided to make a more useful hand. Surgeons did this at the Institute of Reconstructive Plastic Surgery in New York to a 32-year-old man with the initials O. G.. They touched O. G.'s fingers before and after surgery while using MRI brain scans. Before the surgery, the fingers mapped onto his brain were fused close together; afterward, the maps of his individual fingers did indeed separate and take the layout corresponding to a normal hand. Clinical significance Anomalies, injuries and diseases A rare anatomical variation affects 1 in 500 humans, in which the individual has more than the usual number of digits; this is known as polydactyly. A human may also be born without one or more fingers or underdevelopment of some fingers such as symbrachydactyly. Extra fingers can be functional. One individual with seven fingers not only used them but claimed that they "gave him some advantages in playing the piano". Phalanges are commonly fractured. A damaged tendon can cause significant loss of function in fine motor control, such as with a mallet finger. They can be damaged by cold, including frostbite and non-freezing cold injury (NFCI); and heat, including burns. The fingers are commonly affected by diseases such as rheumatoid arthritis and gout. Individuals with diabetes often use the fingers to obtain blood samples for regular blood sugar testing. Raynaud's phenomenon and Paroxysmal hand hematoma are neurovascular disorders that affect the fingers. Research has linked the ratio of lengths between the index and ring fingers to higher levels of testosterone, and to various physical and behavioral traits such as penis length and risk for development of alcohol dependence or video game addiction. Etymology The English word finger stems from Old English finger, ultimately from Proto-Germanic ('finger'). It is cognate with Gothic , Old Norse , or Old High German . Linguists generally assume that is a ro-stem deriving from a previous form , ultimately from Proto-Indo-European ('five'). The name pinkie derives from Dutch , of uncertain origin. In English only the digits on the hand are known as fingers. However, in some languages the translated version of fingers can mean either the digits on the hand or feet. In English a digit on a foot has the distinct name of toe.	Biology and health sciences	Human anatomy	null
37636	https://en.wikipedia.org/wiki/Encephalitis	Encephalitis	Encephalitis is inflammation of the brain. The severity can be variable with symptoms including reduction or alteration in consciousness, aphasia, headache, fever, confusion, a stiff neck, and vomiting. Complications may include seizures, hallucinations, trouble speaking, memory problems, and problems with hearing. Causes of encephalitis include viruses such as herpes simplex virus and rabies virus as well as bacteria, fungi, or parasites. Other causes include autoimmune diseases and certain medications. In many cases the cause remains unknown. Risk factors include a weak immune system. Diagnosis is typically based on symptoms and supported by blood tests, medical imaging, and analysis of cerebrospinal fluid. Certain types are preventable with vaccines. Treatment may include antiviral medications (such as acyclovir), anticonvulsants, and corticosteroids. Treatment generally takes place in hospital. Some people require artificial respiration. Once the immediate problem is under control, rehabilitation may be required. In 2015, encephalitis was estimated to have affected 4.3 million people and resulted in 150,000 deaths worldwide. Signs and symptoms Adults with encephalitis present with acute onset of fever, headache, confusion, and sometimes seizures. Younger children or infants may present with irritability, poor appetite and fever. Neurological examinations usually reveal a drowsy or confused person. Stiff neck, due to the irritation of the meninges covering the brain, indicates that the patient has either meningitis or meningoencephalitis. Limbic encephalitis Limbic encephalitis refers to inflammatory disease confined to the limbic system of the brain. The clinical presentation often includes disorientation, disinhibition, memory loss, seizures, and behavioral anomalies. MRI imaging reveals T2 hyperintensity in the structures of the medial temporal lobes, and in some cases, other limbic structures. Some cases of limbic encephalitis are of autoimmune origin. Encephalitis lethargica Encephalitis lethargica is identified by high fever, headache, delayed physical response, and lethargy. Individuals can exhibit upper body weakness, muscular pains, and tremors, though the cause of encephalitis lethargica is not currently known. From 1917 to 1928, an epidemic of encephalitis lethargica occurred worldwide. Cause In 30%-40% of encephalitis cases, the etiology remains unknown. Viral Viral infections are the usual cause of infectious encephalitis. Viral encephalitis can occur either as a direct effect of an acute infection, or as one of the sequelae of a latent infection. The majority of viral cases of encephalitis have an unknown cause; however, the most common identifiable cause of viral encephalitis is from herpes simplex infection. Other causes of acute viral encephalitis are rabies virus, poliovirus, and measles virus. Additional possible viral causes are arboviral flavivirus (St. Louis encephalitis, West Nile virus), bunyavirus (La Crosse strain), arenavirus (lymphocytic choriomeningitis virus), reovirus (Colorado tick virus), and henipavirus infections. The Powassan virus is a rare cause of encephalitis. Bacterial It can be caused by a bacterial infection, such as bacterial meningitis, or may be a complication of a current infectious disease such as syphilis (secondary encephalitis). Other bacterial pathogens, like Mycoplasma and those causing rickettsial disease, cause inflammation of the meninges and consequently encephalitis. Lyme disease or Bartonella henselae may also cause encephalitis. Other infectious causes Certain parasitic or protozoal infestations, such as toxoplasmosis and malaria can also cause encephalitis in people with compromised immune systems. The rare but typically deadly forms of encephalitis, primary amoebic meningoencephalitis and Granulomatous amoebic encephalitis, are caused by free-living amoeba. Autoimmune encephalitis Autoimmune encephalitis signs can include catatonia, psychosis, abnormal movements, and autonomic dysregulation. Antibody-mediated anti-N-methyl-D-aspartate-receptor encephalitis and Rasmussen encephalitis are examples of autoimmune encephalitis. Anti-NMDA receptor encephalitis is the most common autoimmune form, and is accompanied by ovarian teratoma in 58 percent of affected women 18–45 years of age. Another autoimmune cause includes acute disseminated encephalitis, a demyelinating disease which primarily affects children. Diagnosis People should only be diagnosed with encephalitis if they have a decreased or altered level of consciousness, lethargy, or personality change for at least twenty-four hours without any other explainable cause. Diagnosing encephalitis is done via a variety of tests: Brain scan, done by MRI, can determine inflammation and differentiate from other possible causes. EEG, in monitoring brain activity, encephalitis will produce abnormal signal. Lumbar puncture (spinal tap), this helps determine via a test using the cerebral-spinal fluid, obtained from the lumbar region. Blood test Urine analysis Polymerase chain reaction (PCR) testing of the cerebrospinal fluid, to detect the presence of viral DNA which is a sign of viral encephalitis. Prevention Vaccination is available against tick-borne and Japanese encephalitis and should be considered for at-risk individuals. Post-infectious encephalomyelitis complicating smallpox vaccination is avoidable, for all intents and purposes, as smallpox is nearly eradicated. Contraindication to Pertussis immunization should be observed in patients with encephalitis. Treatment An ideal drug to treat brain infection should be small, moderately lipophilic at pH of 7.4, low level of plasma protein binding, volume of distribution of litre per kg, does not have strong affinity towards binding with P-glycoprotein, or other efflux pumps on the surface of blood–brain barrier. Some drugs such as isoniazid, pyrazinamide, linezolid, metronidazole, fluconazole, and some fluoroquinolones have good penetration to blood brain barrier. Treatment (which is based on supportive care) is as follows: Pyrimethamine-based maintenance therapy is often used to treat toxoplasmic encephalitis (TE), which is caused by Toxoplasma gondii and can be life-threatening for people with weak immune systems. The use of highly active antiretroviral therapy (HAART), in conjunction with the established pyrimethamine-based maintenance therapy, decreases the chance of relapse in patients with HIV and TE from approximately 18% to 11%. This is a significant difference as relapse may impact the severity and prognosis of disease and result in an increase in healthcare expenditure. The effectiveness of intravenous immunoglobulin for the management of childhood encephalitis is unclear. Systematic reviews have been unable to draw firm conclusions because of a lack of randomised double-blind studies with sufficient numbers of patients and sufficient follow-up. There is the possibility of a benefit of intravenous immunoglobulin for some forms of childhood encephalitis on some indicators such as length of hospital stay, time to stop spasms, time to regain consciousness, and time to resolution of neuropathic symptoms and fever. Intravenous immunoglobulin for Japanese encephalitis appeared to have no benefit when compared with placebo (pretend) treatment. Prognosis Identification of poor prognostic factors include cerebral edema, status epilepticus, and thrombocytopenia. In contrast, a normal encephalogram at the early stages of diagnosis is associated with high rates of survival. Epidemiology The number of new cases a year of acute encephalitis in Western countries is 7.4 cases per 100,000 people per year. In tropical countries, the incidence is 6.34 per 100,000 people per year. The number of cases of encephalitis has not changed much over time, with about 250,000 cases a year from 2005 to 2015 in the US. Approximately seven per 100,000 people were hospitalized for encephalitis in the US during this time. In 2015, encephalitis was estimated to have affected 4.3 million people and resulted in 150,000 deaths worldwide. Herpes simplex encephalitis has an incidence of 2–4 per million of the population per year. Terminology Encephalitis with meningitis is known as meningoencephalitis, while encephalitis with involvement of the spinal cord is known as encephalomyelitis. The word is from Ancient Greek , 'brain', composed of , , 'in' and , , 'head', and the medical suffix -itis 'inflammation'.	Biology and health sciences	Infectious disease	null
37637	https://en.wikipedia.org/wiki/Nucleophile	Nucleophile	In chemistry, a nucleophile is a chemical species that forms bonds by donating an electron pair. All molecules and ions with a free pair of electrons or at least one pi bond can act as nucleophiles. Because nucleophiles donate electrons, they are Lewis bases. Nucleophilic describes the affinity of a nucleophile to bond with positively charged atomic nuclei. Nucleophilicity, sometimes referred to as nucleophile strength, refers to a substance's nucleophilic character and is often used to compare the affinity of atoms. Neutral nucleophilic reactions with solvents such as alcohols and water are named solvolysis. Nucleophiles may take part in nucleophilic substitution, whereby a nucleophile becomes attracted to a full or partial positive charge, and nucleophilic addition. Nucleophilicity is closely related to basicity. The difference between the two is, that basicity is a thermodynamic property (i.e. relates to an equilibrium state), but nucleophilicity is a kinetic property, which relates to rates of certain chemical reactions. History and Etymology The terms nucleophile and electrophile were introduced by Christopher Kelk Ingold in 1933, replacing the terms anionoid and cationoid proposed earlier by A. J. Lapworth in 1925. The word nucleophile is derived from nucleus and the Greek word φιλος, philos, meaning friend. Properties In general, in a group across the periodic table, the more basic the ion (the higher the pKa of the conjugate acid) the more reactive it is as a nucleophile. Within a series of nucleophiles with the same attacking element (e.g. oxygen), the order of nucleophilicity will follow basicity. Sulfur is in general a better nucleophile than oxygen. Nucleophilicity Many schemes attempting to quantify relative nucleophilic strength have been devised. The following empirical data have been obtained by measuring reaction rates for many reactions involving many nucleophiles and electrophiles. Nucleophiles displaying the so-called alpha effect are usually omitted in this type of treatment. Swain–Scott equation The first such attempt is found in the Swain–Scott equation derived in 1953: This free-energy relationship relates the pseudo first order reaction rate constant (in water at 25 °C), k, of a reaction, normalized to the reaction rate, k0, of a standard reaction with water as the nucleophile, to a nucleophilic constant n for a given nucleophile and a substrate constant s that depends on the sensitivity of a substrate to nucleophilic attack (defined as 1 for methyl bromide). This treatment results in the following values for typical nucleophilic anions: acetate 2.7, chloride 3.0, azide 4.0, hydroxide 4.2, aniline 4.5, iodide 5.0, and thiosulfate 6.4. Typical substrate constants are 0.66 for ethyl tosylate, 0.77 for β-propiolactone, 1.00 for 2,3-epoxypropanol, 0.87 for benzyl chloride, and 1.43 for benzoyl chloride. The equation predicts that, in a nucleophilic displacement on benzyl chloride, the azide anion reacts 3000 times faster than water. Ritchie equation The Ritchie equation, derived in 1972, is another free-energy relationship: where N+ is the nucleophile dependent parameter and k0 the reaction rate constant for water. In this equation, a substrate-dependent parameter like s in the Swain–Scott equation is absent. The equation states that two nucleophiles react with the same relative reactivity regardless of the nature of the electrophile, which is in violation of the reactivity–selectivity principle. For this reason, this equation is also called the constant selectivity relationship. In the original publication the data were obtained by reactions of selected nucleophiles with selected electrophilic carbocations such as tropylium or diazonium cations: or (not displayed) ions based on malachite green. Many other reaction types have since been described. Typical Ritchie N+ values (in methanol) are: 0.5 for methanol, 5.9 for the cyanide anion, 7.5 for the methoxide anion, 8.5 for the azide anion, and 10.7 for the thiophenol anion. The values for the relative cation reactivities are −0.4 for the malachite green cation, +2.6 for the benzenediazonium cation, and +4.5 for the tropylium cation. Mayr–Patz equation In the Mayr–Patz equation (1994): The second order reaction rate constant k at 20 °C for a reaction is related to a nucleophilicity parameter N, an electrophilicity parameter E, and a nucleophile-dependent slope parameter s. The constant s is defined as 1 with 2-methyl-1-pentene as the nucleophile. Many of the constants have been derived from reaction of so-called benzhydrylium ions as the electrophiles: and a diverse collection of π-nucleophiles: . Typical E values are +6.2 for R = chlorine, +5.90 for R = hydrogen, 0 for R = methoxy and −7.02 for R = dimethylamine. Typical N values with s in parentheses are −4.47 (1.32) for electrophilic aromatic substitution to toluene (1), −0.41 (1.12) for electrophilic addition to 1-phenyl-2-propene (2), and 0.96 (1) for addition to 2-methyl-1-pentene (3), −0.13 (1.21) for reaction with triphenylallylsilane (4), 3.61 (1.11) for reaction with 2-methylfuran (5), +7.48 (0.89) for reaction with isobutenyltributylstannane (6) and +13.36 (0.81) for reaction with the enamine 7. The range of organic reactions also include SN2 reactions: With E = −9.15 for the S-methyldibenzothiophenium ion, typical nucleophile values N (s) are 15.63 (0.64) for piperidine, 10.49 (0.68) for methoxide, and 5.20 (0.89) for water. In short, nucleophilicities towards sp2 or sp3 centers follow the same pattern. Unified equation In an effort to unify the above described equations the Mayr equation is rewritten as: with sE the electrophile-dependent slope parameter and sN the nucleophile-dependent slope parameter. This equation can be rewritten in several ways: with sE = 1 for carbocations this equation is equal to the original Mayr–Patz equation of 1994, with sN = 0.6 for most n nucleophiles the equation becomes or the original Scott–Swain equation written as: with sE = 1 for carbocations and sN = 0.6 the equation becomes: or the original Ritchie equation written as: Types Examples of nucleophiles are anions such as Cl−, or a compound with a lone pair of electrons such as NH3 (ammonia) and PR3. In the example below, the oxygen of the hydroxide ion donates an electron pair to form a new chemical bond with the carbon at the end of the bromopropane molecule. The bond between the carbon and the bromine then undergoes heterolytic fission, with the bromine atom taking the donated electron and becoming the bromide ion (Br−), because a SN2 reaction occurs by backside attack. This means that the hydroxide ion attacks the carbon atom from the other side, exactly opposite the bromine ion. Because of this backside attack, SN2 reactions result in a inversion of the configuration of the electrophile. If the electrophile is chiral, it typically maintains its chirality, though the SN2 product's absolute configuration is flipped as compared to that of the original electrophile. Ambident Nucleophile An ambident nucleophile is one that can attack from two or more places, resulting in two or more products. For example, the thiocyanate ion (SCN−) may attack from either the sulfur or the nitrogen. For this reason, the SN2 reaction of an alkyl halide with SCN− often leads to a mixture of an alkyl thiocyanate (R-SCN) and an alkyl isothiocyanate (R-NCS). Similar considerations apply in the Kolbe nitrile synthesis. Halogens While the halogens are not nucleophilic in their diatomic form (e.g. I2 is not a nucleophile), their anions are good nucleophiles. In polar, protic solvents, F− is the weakest nucleophile, and I− the strongest; this order is reversed in polar, aprotic solvents. Carbon Carbon nucleophiles are often organometallic reagents such as those found in the Grignard reaction, Blaise reaction, Reformatsky reaction, and Barbier reaction or reactions involving organolithium reagents and acetylides. These reagents are often used to perform nucleophilic additions. Enols are also carbon nucleophiles. The formation of an enol is catalyzed by acid or base. Enols are ambident nucleophiles, but, in general, nucleophilic at the alpha carbon atom. Enols are commonly used in condensation reactions, including the Claisen condensation and the aldol condensation reactions. Oxygen Examples of oxygen nucleophiles are water (H2O), hydroxide anion, alcohols, alkoxide anions, hydrogen peroxide, and carboxylate anions. Nucleophilic attack does not take place during intermolecular hydrogen bonding. Sulfur Of sulfur nucleophiles, hydrogen sulfide and its salts, thiols (RSH), thiolate anions (RS−), anions of thiolcarboxylic acids (RC(O)-S−), and anions of dithiocarbonates (RO-C(S)-S−) and dithiocarbamates (R2N-C(S)-S−) are used most often. In general, sulfur is very nucleophilic because of its large size, which makes it readily polarizable, and its lone pairs of electrons are readily accessible. Nitrogen Nitrogen nucleophiles include ammonia, azide, amines, nitrites, hydroxylamine, hydrazine, carbazide, phenylhydrazine, semicarbazide, and amide. Metal centers Although metal centers (e.g., Li+, Zn2+, Sc3+, etc.) are most commonly cationic and electrophilic (Lewis acidic) in nature, certain metal centers (particularly ones in a low oxidation state and/or carrying a negative charge) are among the strongest recorded nucleophiles and are sometimes referred to as "supernucleophiles." For instance, using methyl iodide as the reference electrophile, Ph3Sn– is about 10000 times more nucleophilic than I–, while the Co(I) form of vitamin B12 (vitamin B12s) is about 107 times more nucleophilic. Other supernucleophilic metal centers include low oxidation state carbonyl metalate anions (e.g., CpFe(CO)2–). Examples The following table shows the nucleophilicity of some molecules with methanol as the solvent:	Physical sciences	Concepts	Chemistry
37649	https://en.wikipedia.org/wiki/Amplitude	Amplitude	The amplitude of a periodic variable is a measure of its change in a single period (such as time or spatial period). The amplitude of a non-periodic signal is its magnitude compared with a reference value. There are various definitions of amplitude (see below), which are all functions of the magnitude of the differences between the variable's extreme values. In older texts, the phase of a periodic function is sometimes called the amplitude. Definitions Peak amplitude and semi-amplitude For symmetric periodic waves, like sine waves or triangle waves, peak amplitude and semi amplitude are the same. Peak amplitude In audio system measurements, telecommunications and others where the measurand is a signal that swings above and below a reference value but is not sinusoidal, peak amplitude is often used. If the reference is zero, this is the maximum absolute value of the signal; if the reference is a mean value (DC component), the peak amplitude is the maximum absolute value of the difference from that reference. Semi-amplitude Semi-amplitude means half of the peak-to-peak amplitude. The majority of scientific literature employs the term amplitude or peak amplitude to mean semi-amplitude. It is the most widely used measure of orbital wobble in astronomy and the measurement of small radial velocity semi-amplitudes of nearby stars is important in the search for exoplanets (see Doppler spectroscopy). Ambiguity In general, the use of peak amplitude is simple and unambiguous only for symmetric periodic waves, like a sine wave, a square wave, or a triangle wave. For an asymmetric wave (periodic pulses in one direction, for example), the peak amplitude becomes ambiguous. This is because the value is different depending on whether the maximum positive signal is measured relative to the mean, the maximum negative signal is measured relative to the mean, or the maximum positive signal is measured relative to the maximum negative signal (the peak-to-peak amplitude) and then divided by two (the semi-amplitude). In electrical engineering, the usual solution to this ambiguity is to measure the amplitude from a defined reference potential (such as ground or 0 V). Strictly speaking, this is no longer amplitude since there is the possibility that a constant (DC component) is included in the measurement. Peak-to-peak amplitude Peak-to-peak amplitude (abbreviated p–p or PtP or PtoP) is the change between peak (highest amplitude value) and trough (lowest amplitude value, which can be negative). With appropriate circuitry, peak-to-peak amplitudes of electric oscillations can be measured by meters or by viewing the waveform on an oscilloscope. Peak-to-peak is a straightforward measurement on an oscilloscope, the peaks of the waveform being easily identified and measured against the graticule. This remains a common way of specifying amplitude, but sometimes other measures of amplitude are more appropriate. Root mean square amplitude Root mean square (RMS) amplitude is used especially in electrical engineering: the RMS is defined as the square root of the mean over time of the square of the vertical distance of the graph from the rest state; i.e. the RMS of the AC waveform (with no DC component). For complicated waveforms, especially non-repeating signals like noise, the RMS amplitude is usually used because it is both unambiguous and has physical significance. For example, the average power transmitted by an acoustic or electromagnetic wave or by an electrical signal is proportional to the square of the RMS amplitude (and not, in general, to the square of the peak amplitude). For alternating current electric power, the universal practice is to specify RMS values of a sinusoidal waveform. One property of root mean square voltages and currents is that they produce the same heating effect as a direct current in a given resistance. The peak-to-peak value is used, for example, when choosing rectifiers for power supplies, or when estimating the maximum voltage that insulation must withstand. Some common voltmeters are calibrated for RMS amplitude, but respond to the average value of a rectified waveform. Many digital voltmeters and all moving coil meters are in this category. The RMS calibration is only correct for a sine wave input since the ratio between peak, average and RMS values is dependent on waveform. If the wave shape being measured is greatly different from a sine wave, the relationship between RMS and average value changes. True RMS-responding meters were used in radio frequency measurements, where instruments measured the heating effect in a resistor to measure a current. The advent of microprocessor-controlled meters capable of calculating RMS by sampling the waveform has made true RMS measurement commonplace. Pulse amplitude In telecommunications, pulse amplitude is the magnitude of a pulse parameter, such as the voltage level, current level, field intensity, or power level. Pulse amplitude is measured with respect to a specified reference and therefore should be modified by qualifiers, such as average, instantaneous, peak, or root-mean-square. Pulse amplitude also applies to the amplitude of frequency- and phase-modulated waveform envelopes. Formal representation In this simple wave equation is the amplitude (or peak amplitude), is the oscillating variable, is angular frequency, is time, and are arbitrary constants representing time and displacement offsets respectively. Units The units of the amplitude depend on the type of wave, but are always in the same units as the oscillating variable. A more general representation of the wave equation is more complex, but the role of amplitude remains analogous to this simple case. For waves on a string, or in a medium such as water, the amplitude is a displacement. The amplitude of sound waves and audio signals (which relates to the volume) conventionally refers to the amplitude of the air pressure in the wave, but sometimes the amplitude of the displacement (movements of the air or the diaphragm of a speaker) is described. The logarithm of the amplitude squared is usually quoted in dB, so a null amplitude corresponds to −∞ dB. Loudness is related to amplitude and intensity and is one of the most salient qualities of a sound, although in general sounds it can be recognized independently of amplitude. The square of the amplitude is proportional to the intensity of the wave. For electromagnetic radiation, the amplitude of a photon corresponds to the changes in the electric field of the wave. However, radio signals may be carried by electromagnetic radiation; the intensity of the radiation (amplitude modulation) or the frequency of the radiation (frequency modulation) is oscillated and then the individual oscillations are varied (modulated) to produce the signal. Amplitude envelopes Amplitude envelope refers to the changes in the amplitude of a sound over time, and is an influential property as it affects perception of timbre. A flat tone has a steady state amplitude that remains constant during time, which is represented by a scalar. Other sounds can have percussive amplitude envelopes featuring an abrupt onset followed by an immediate exponential decay. Percussive amplitude envelopes are characteristic of various impact sounds: two wine glasses clinking together, hitting a drum, slamming a door, etc. where the amplitude is transient and must be represented as either a continuous function or a discrete vector. Percussive amplitude envelopes model many common sounds that have a transient loudness attack, decay, sustain, and release. Amplitude normalization With waveforms containing many overtones, complex transient timbres can be achieved by assigning each overtone to its own distinct transient amplitude envelope. Unfortunately, this has the effect of modulating the loudness of the sound as well. It makes more sense to separate loudness and harmonic quality to be parameters controlled independently of each other. To do so, harmonic amplitude envelopes are frame-by-frame normalized to become amplitude proportion envelopes, where at each time frame all the harmonic amplitudes will add to 100% (or 1). This way, the main loudness-controlling envelope can be cleanly controlled. In Sound Recognition, max amplitude normalization can be used to help align the key harmonic features of 2 alike sounds, allowing similar timbres to be recognized independent of loudness.	Physical sciences	Waves	null
37654	https://en.wikipedia.org/wiki/Owl	Owl	Owls are birds from the order Strigiformes (), which includes over 200 species of mostly solitary and nocturnal birds of prey typified by an upright stance, a large, broad head, binocular vision, binaural hearing, sharp talons, and feathers adapted for silent flight. Exceptions include the diurnal northern hawk-owl and the gregarious burrowing owl. Owls are divided into two families: the true (or typical) owl family, Strigidae, and the barn owl and bay owl family, Tytonidae. Owls hunt mostly small mammals, insects, and other birds, although a few species specialize in hunting fish. They are found in all regions of the Earth except the polar ice caps and some remote islands. A group of owls is called a "parliament". Anatomy Owls possess large, forward-facing eyes and ear-holes, a hawk-like beak, a flat face, and usually a conspicuous circle of feathers, a facial disc, around each eye. The feathers making up this disc can be adjusted to sharply focus sounds from varying distances onto the owls' asymmetrically placed ear cavities. Most birds of prey have eyes on the sides of their heads, but the stereoscopic nature of the owl's forward-facing eyes permits the greater sense of depth perception necessary for low-light hunting. Owls have binocular vision, but they must rotate their entire heads to change the focus of their view because, like most birds, their eyes are fixed in their sockets. Owls are farsighted and cannot clearly see anything nearer than a few centimetres of their eyes. Caught prey can be felt by owls with the use of filoplumes—hairlike feathers on the beak and feet that act as "feelers". Their far vision, particularly in low light, is exceptionally good. Owls can rotate their heads and necks as much as 270°. Owls have 14 neck vertebrae — humans have only seven — and their vertebral circulatory systems are adapted to allow them to rotate their heads without cutting off blood to the brain. Specifically, the foramina in their vertebrae through which the vertebral arteries pass are about ten times the diameter of the artery, instead of about the same size as the artery, as is the case in humans; the vertebral arteries enter the cervical vertebrae higher than in other birds, giving the vessels some slack, and the carotid arteries unite in a very large anastomosis or junction, the largest of any bird's, preventing blood supply from being cut off while they rotate their necks. Other anastomoses between the carotid and vertebral arteries support this effect. The smallest owl—weighing as little as and measuring some —is the elf owl (Micrathene whitneyi). Around the same diminutive length, although slightly heavier, are the lesser known long-whiskered owlet (Xenoglaux loweryi) and Tamaulipas pygmy owl (Glaucidium sanchezi). The largest owls are two similarly sized eagle owls; the Eurasian eagle-owl (Bubo bubo) and Blakiston's fish owl (Bubo blakistoni). The largest females of these species are long, have a wing span, and weigh . Different species of owls produce different sounds; this distribution of calls aids owls in finding mates or announcing their presence to potential competitors, and also aids ornithologists and birders in locating these birds and distinguishing species. As noted above, their facial discs help owls to funnel the sound of prey to their ears. In many species, these discs are placed asymmetrically, for better directional location. Owl plumage is generally cryptic, although several species have facial and head markings, including face masks, ear tufts, and brightly colored irises. These markings are generally more common in species inhabiting open habitats, and are thought to be used in signaling with other owls in low-light conditions. Sexual dimorphism Sexual dimorphism is a physical difference between males and females of a species. Female owls are typically larger than the males. The degree of size dimorphism varies across multiple populations and species, and is measured through various traits, such as wing span and body mass. One theory suggests that selection has led males to be smaller because it allows them to be efficient foragers. The ability to obtain more food is advantageous during breeding season. In some species, female owls stay at their nest with their eggs while it is the responsibility of the male to bring back food to the nest. If food is scarce, the male first feeds himself before feeding the female. Small birds, which are agile, are an important source of food for owls. Male burrowing owls have been observed to have longer wing chords than females, despite being smaller than females. Furthermore, owls have been observed to be roughly the same size as their prey. This has also been observed in other predatory birds, which suggests that owls with smaller bodies and long wing chords have been selected for because of the increased agility and speed that allows them to catch their prey. Another popular theory suggests that females have not been selected to be smaller like male owls because of their sexual roles. In many species, female owls may not leave the nest. Therefore, females may have a larger mass to allow them to go for a longer period of time without starving. For example, one hypothesized sexual role is that larger females are more capable of dismembering prey and feeding it to their young, hence female owls are larger than their male counterparts. A different theory suggests that the size difference between male and females is due to sexual selection: since large females can choose their mate and may violently reject a male's sexual advances, smaller male owls that have the ability to escape unreceptive females are more likely to have been selected. If the character is stable, there can be different optimums for both sexes. Selection operates on both sexes at the same time; therefore it is necessary to explain not only why one of the sexes is relatively bigger, but also why the other sex is smaller. If owls are still evolving toward smaller bodies and longer wing chords, according to V. Geodakyan's Evolutionary Theory of Sex, males should be more advanced on these characters. Males are viewed as an evolutionary vanguard of a population, and sexual dimorphism on the character, as an evolutionary "distance" between the sexes. "Phylogenetic rule of sexual dimorphism" states that if there exists a sexual dimorphism on any character, then the evolution of this trait goes from the female form toward the male one. Hunting adaptations All owls are carnivorous birds of prey and live on diets of insects, small rodents and lagomorphs. Some owls are also specifically adapted to hunt fish. They are very adept in hunting in their respective environments. Since owls can be found in nearly all parts of the world and across a multitude of ecosystems, their hunting skills and characteristics vary slightly from species to species, though most characteristics are shared among all species. Flight and feathers Most owls share an innate ability to fly almost silently and also more slowly in comparison to other birds of prey. Most owls live a mainly nocturnal lifestyle and being able to fly without making any noise gives them a strong advantage over prey alert to the slightest sound in the night. A silent, slow flight is not as necessary for diurnal and crepuscular owls given that prey can usually see an owl approaching. Owls' feathers are generally larger than the average birds' feathers, have fewer radiates, longer pennulum, and achieve smooth edges with different rachis structures. Serrated edges along the owl's remiges bring the flapping of the wing down to a nearly silent mechanism. The serrations are more likely reducing aerodynamic disturbances, rather than simply reducing noise. The surface of the flight feathers is covered with a velvety structure that absorbs the sound of the wing moving. These unique structures reduce noise frequencies above 2 kHz, making the sound level emitted drop below the typical hearing spectrum of the owl's usual prey and also within the owl's own best hearing range. This optimizes the owl's ability to silently fly to capture prey without the prey hearing the owl first as it flies, and to hear any noise the prey makes. It also allows the owl to monitor the sound output from its flight pattern. The disadvantage of such feather adaptations for barn owls is that their feathers are not waterproof. The adaptations mean that barn owls do not use the uropygial gland, informally the "preen" or "oil" gland, as most birds do, to spread oils across their plumage through preening. This makes them highly vulnerable to heavy rain when they are unable to hunt. Historically, they would switch to hunting indoors in wet weather, using barns and other agricultural buildings, but the decline in the numbers of these structures in the 20th and 21st centuries has reduced such opportunities. The lack of waterproofing means that barn owls are also susceptible to drowning, in drinking troughs and other structures with smooth sides. The Barn Owl Trust provides advice on how this can be mitigated, by the installation of floats. Vision Eyesight is a particular characteristic of the owl that aids in nocturnal prey capture. Owls are part of a small group of birds that live nocturnally, but do not use echolocation to guide them in flight in low-light situations. Owls are known for their disproportionally large eyes in comparison to their skulls. An apparent consequence of the evolution of an absolutely large eye in a relatively small skull is that the eye of the owl has become tubular in shape. This shape is found in other so-called nocturnal eyes, such as the eyes of strepsirrhine primates and bathypelagic fishes. Since the eyes are fixed into these sclerotic tubes, they are unable to move the eyes in any direction. Instead of moving their eyes, owls swivel their heads to view their surroundings. Owls' heads are capable of swiveling through an angle of roughly 270°, easily enabling them to see behind them without relocating the torso. This ability keeps bodily movement at a minimum, thus reduces the amount of sound the owl makes as it waits for its prey. Owls are regarded as having the most frontally placed eyes among all avian groups, which gives them some of the largest binocular fields of vision. Owls are farsighted and cannot focus on objects within a few centimetres of their eyes. These mechanisms are only able to function due to the large-sized retinal image. Thus, the primary nocturnal function in the vision of the owl is due to its large posterior nodal distance; retinal image brightness is only maximized to the owl within secondary neural functions. These attributes of the owl cause its nocturnal eyesight to be far superior to that of its average prey. Hearing Owls exhibit specialized hearing functions and ear shapes that also aid in hunting. They are noted for asymmetrical ear placements on the skull in some genera. Owls can have either internal or external ears, both of which are asymmetrical. Asymmetry has not been reported to extend to the middle or internal ear of the owl. Asymmetrical ear placement on the skull allows the owl to pinpoint the location of its prey. This is especially true for strictly nocturnal species such as the barn owls Tyto or Tengmalm's owl. With ears set at different places on its skull, an owl is able to determine the direction from which the sound is coming by the minute difference in time that it takes for the sound waves to penetrate the left and right ears. The owl turns its head until the sound reaches both ears at the same time, at which point it is directly facing the source of the sound. This time difference between ears is about 30 microseconds. Behind the ear openings are modified, dense feathers, densely packed to form a facial ruff, which creates an anterior-facing, concave wall that cups the sound into the ear structure. This facial ruff is poorly defined in some species, and prominent, nearly encircling the face, in other species. The facial disk also acts to direct sound into the ears, and a downward-facing, sharply triangular beak minimizes sound reflection away from the face. The shape of the facial disk is adjustable at will to focus sounds more effectively. The prominences above a great horned owl's head are commonly mistaken as its ears. This is not the case; they are merely feather tufts. The ears are on the sides of the head in the usual location (in two different locations as described above). Talons While the auditory and visual capabilities of the owl allow it to locate and pursue its prey, the talons and beak of the owl do the final work. The owl kills its prey using these talons to crush the skull and knead the body. The crushing power of an owl's talons varies according to prey size and type, and by the size of the owl. The burrowing owl (Athene cunicularia), a small, partly insectivorous owl, has a release force of only 5 N. The larger barn owl (Tyto alba) needs a force of 30 N to release its prey, and one of the largest owls, the great horned owl (Bubo virginianus), needs a force over 130 N to release prey in its talons. An owl's talons, like those of most birds of prey, can seem massive in comparison to the body size outside of flight. The Tasmanian masked owl has some of the proportionally longest talons of any bird of prey; they appear enormous in comparison to the body when fully extended to grasp prey. An owl's claws are sharp and curved. The family Tytonidae has inner and central toes of about equal length, while the family Strigidae has an inner toe that is distinctly shorter than the central one. These different morphologies allow efficiency in capturing prey specific to the different environments they inhabit. Beak The beak of the owl is short, curved, and downward-facing, and typically hooked at the tip for gripping and tearing its prey. Once prey is captured, the scissor motion of the top and lower bill is used to tear the tissue and kill. The sharp lower edge of the upper bill works in coordination with the sharp upper edge of the lower bill to deliver this motion. The downward-facing beak allows the owl's field of vision to be clear, as well as directing sound into the ears without deflecting sound waves away from the face. Camouflage The coloration of the owl's plumage plays a key role in its ability to sit still and blend into the environment, making it nearly invisible to prey. Owls tend to mimic the coloration and sometimes the texture patterns of their surroundings, the barn owl being an exception. The snowy owl (Bubo scandiacus) appears nearly bleach-white in color with a few flecks of black, mimicking their snowy surroundings perfectly, while the speckled brown plumage of the tawny owl (Strix aluco) allows it to lie in wait among the deciduous woodland it prefers for its habitat. Likewise, the mottled wood owl (Strix ocellata) displays shades of brown, tan, and black, making the owl nearly invisible in the surrounding trees, especially from behind. Usually, the only tell-tale sign of a perched owl is its vocalizations or its vividly colored eyes. Behavior Most owls are nocturnal, actively hunting their prey in darkness. Several types of owls are crepuscular—active during the twilight hours of dawn and dusk; one example is the pygmy owl (Glaucidium). A few owls are active during the day, also; examples are the burrowing owl (Speotyto cunicularia) and the short-eared owl (Asio flammeus). Much of the owls' hunting strategy depends on stealth and surprise. Owls have at least two adaptations that aid them in achieving stealth. First, the dull coloration of their feathers can render them almost invisible under certain conditions. Secondly, serrated edges on the leading edge of owls' remiges muffle an owl's wing beats, allowing an owl's flight to be practically silent. Some fish-eating owls, for which silence has no evolutionary advantage, lack this adaptation. An owl's sharp beak and powerful talons allow it to kill its prey before swallowing it whole (if it is not too big). Scientists studying the diets of owls are helped by their habit of regurgitating the indigestible parts of their prey (such as bones, scales, and fur) in the form of pellets. These "owl pellets" are plentiful and easy to interpret, and are often sold by companies to schools for dissection by students as a lesson in biology and ecology. Breeding and reproduction Owl eggs typically have a white color and an almost spherical shape, and range in number from a few to a dozen, depending on species and the particular season; for most, three or four is the more common number. In at least one species, female owls do not mate with the same male for a lifetime. Female burrowing owls commonly travel and find other mates, while the male stays in his territory and mates with other females. Evolution and systematics Recent phylogenetic studies place owls within the clade Telluraves, most closely related to the Accipitrimorphae and the Coraciimorphae, although the exact placement within Telluraves is disputed. See below cladogram: Cladogram of Telluraves relationships based on Braun & Kimball (2021) Some 220 to 225 extant species of owls are known, subdivided into two families: 1. true owls or typical owls family (Strigidae) and 2. barn-owls family (Tytonidae). Some entirely extinct families have also been erected based on fossil remains; these differ much from modern owls in being less specialized or specialized in a very different way (such as the terrestrial Sophiornithidae). The Paleocene genera Berruornis and Ogygoptynx show that owls were already present as a distinct lineage some 60–57 million years ago (Mya), hence, possibly also some 5 million years earlier, at the extinction of the non-avian dinosaurs. This makes them one of the oldest known groups of non-Galloanserae landbirds. The supposed "Cretaceous owls" Bradycneme and Heptasteornis are apparently non-avialan maniraptors. During the Paleogene, the Strigiformes radiated into ecological niches now mostly filled by other groups of birds. The owls as known today, though, evolved their characteristic morphology and adaptations during that time, too. By the early Neogene, the other lineages had been displaced by other bird orders, leaving only barn owls and typical owls. The latter at that time was usually a fairly generic type of (probably earless) owl similar to today's North American spotted owl or the European tawny owl; the diversity in size and ecology found in typical owls today developed only subsequently. Around the Paleogene-Neogene boundary (some 25 Mya), barn owls were the dominant group of owls in southern Europe and adjacent Asia at least; the distribution of fossil and present-day owl lineages indicates that their decline is contemporary with the evolution of the different major lineages of true owls, which for the most part seems to have taken place in Eurasia. In the Americas, rather, an expansion of immigrant lineages of ancestral typical owls occurred. The supposed fossil herons "Ardea" perplexa (Middle Miocene of Sansan, France) and "Ardea" lignitum (Late Pliocene of Germany) were more probably owls; the latter was apparently close to the modern genus Bubo. Judging from this, the Late Miocene remains from France described as "Ardea" aureliensis should also be restudied. The Messelasturidae, some of which were initially believed to be basal Strigiformes, are now generally accepted to be diurnal birds of prey showing some convergent evolution toward owls. The taxa often united under Strigogyps were formerly placed in part with the owls, specifically the Sophiornithidae; they appear to be Ameghinornithidae instead. For fossil species and paleosubspecies of extant taxa, see the genus and species articles. For a full list of extant and recently extinct owls, see the article List of owl species. Unresolved and basal forms (all fossil) Berruornis (Late Paleocene of France) basal? Sophornithidae? Strigiformes gen. et sp. indet. (Late Paleocene of Zhylga, Kazakhstan) Primoptynx (Early Eocene of Wyoming, U.S.) Palaeoglaux (Middle-Late Eocene of West-Central Europe) own family Palaeoglaucidae or Strigidae? Palaeobyas (Late Eocene/Early Oligocene of Quercy, France) Tytonidae? Sophiornithidae? Palaeotyto (Late Eocene/Early Oligocene of Quercy, France) Tytonidae? Sophiornithidae? Strigiformes gen. et spp. indet. (Early Oligocene of Wyoming, U.S.) Ypresiglaux (Early Eocene of Essex, United Kingdom and Virginia, U.S.) Ogygoptyngidae Ogygoptynx (Middle/Late Paleocene of Colorado, U.S.) Protostrigidae Eostrix (Early Eocene of United States, Europe, and Mongolia). E. gulottai is the smallest known fossil (or living) owl. Minerva (Middle – Late Eocene of western U.S.) formerly Protostrix, includes "Aquila" ferox, "Aquila" lydekkeri, and "Bubo" leptosteus Oligostrix (mid-Oligocene of Saxony, Germany) Sophiornithidae Sophiornis Tytonidae Genus Tyto – the barn owls, grass owls, and masked owls, stand up to tall; some 15 extant species and possibly one recently extinct Genus Phodilus – the bay owls, two to three extant species and possibly one recently extinct Fossil genera Nocturnavis (Late Eocene/Early Oligocene) includes "Bubo" incertus Selenornis (Late Eocene/Early Oligocene) – includes "Asio" henrici Necrobyas (Late Eocene/Early Oligocene – Late Miocene) includes "Bubo" arvernensis and Paratyto Prosybris (Early Oligocene? – Early Miocene) Placement unresolved Tytonidae gen. et sp. indet. "TMT 164" (Middle Miocene) – Prosybris? Strigidae Genus Aegolius – the saw-whet owls, four species Genus Asio – the eared owls, eight species Genus Athene – two to four species (depending on whether the genera Speotyto and Heteroglaux are included or not) Genus Bubo – the horned owls, eagle-owls and fish-owls; paraphyletic with the genera Nyctea, Ketupa, and Scotopelia, some 25 species Genus Glaucidium – the pygmy owls, about 30–35 species Genus Gymnasio – the Puerto Rican owl Genus Gymnoglaux – the bare-legged owl or Cuban screech-owl Genus Lophostrix – the crested owl Genus Jubula – the maned owl Genus Megascops – the screech owls, some 20 species Genus Micrathene – the elf owl Genus Ninox – the Australasian hawk-owls or boobooks, some 20 species Genus Otus – the scops owls; probably paraphyletic, about 45 species Genus Pseudoscops – the Jamaican owl Genus Psiloscops – the flammulated owl Genus Ptilopsis – the white-faced owls, two species Genus Pulsatrix – the spectacled owls, three species Genus Strix – the earless owls, about 15 species, including four previously assigned to Ciccaba Genus Surnia – the northern hawk-owl Genus Taenioptynx - the collared owlet Genus Uroglaux – the Papuan hawk-owl Genus Xenoglaux – the long-whiskered owlet Extinct genera Genus Grallistrix – the stilt-owls, four species; prehistoric Genus Ornimegalonyx – the Caribbean giant owls, one to two species; prehistoric Fossil genera Mioglaux (Late Oligocene? – Early Miocene of West-Central Europe) – includes "Bubo" poirreiri Intutula (Early/Middle – ?Late Miocene of Central Europe) – includes "Strix/Ninox" brevis Alasio (Middle Miocene of Vieux-Collonges, France) – includes "Strix" collongensis Oraristrix – the Brea owl (Late Pleistocene) Placement unresolved "Otus/Strix" wintershofensis: fossil (Early/Middle Miocene of Wintershof West, Germany) – may be close to extant genus Ninox "Strix" edwardsi – fossil (Middle/Late? Miocene) "Asio" pygmaeus – fossil (Early Pliocene of Odesa, Ukraine) Strigidae gen. et sp. indet. UMMP V31030 (Late Pliocene) – Strix/Bubo? the Ibizan owl, Strigidae gen. et sp. indet. – prehistoric Symbolism and mythology Among the Kikuyu of Kenya, it was believed that owls were harbingers of death. If one saw an owl or heard its hoot, someone was going to die. In general, owls are viewed as harbingers of bad luck, ill health, or death. The belief is widespread even today. In Hinduism, an owl is the vahana (mount) of the goddess Lakshmi, especially in the eastern region of India. Owls are considered a symbol of wealth, prosperity, wisdom, good luck, and fortune. This is the reason why Owls are seen with Lakshmi, who is also the goddess of fortune, wealth, and prosperity. At the same time, owls are also associated with evil times in Hinduism. At times, Chamunda (fearsome form of Chandi) is depicted seated on an owl, her vahana (mount or vehicle). Hindus believe that owls are messengers of death. In China, owls were traditionally considered to be omens of evil or misfortune. In Japan, owls are regarded as lucky, although in ancient times they were associated with death. In Sumerian, Akkadian, and Babylonian culture, the owl was associated with Lilith. The modern West generally associates owls with wisdom and vigilance. This link goes back at least as far as Ancient Greece, where Athens, noted for art and scholarship, and Athena, Athens' patron goddess and the goddess of wisdom, had the owl as a symbol. Marija Gimbutas traces veneration of the owl as a goddess, among other birds, to the culture of Old Europe, long pre-dating Indo-European cultures. T. F. Thiselton-Dyer, in his 1883 Folk-lore of Shakespeare, says that from the earliest period it has been considered a bird of ill-omen," and Pliny tells us how, on one occasion, even Rome itself underwent a lustration, because one of them strayed into the Capitol. He represents it also as a funereal bird, a monster of the night, the very abomination of humankind. Virgil describes its death howl from the top of the temple by night, a circumstance introduced as a precursor of Dido's death. Ovid, too, constantly speaks of this bird's presence as an evil omen; and indeed the same notions respecting it may be found among the writings of most of the ancient poets. Native American cultures People often allude to the reputation of owls as bearers of supernatural danger when they tell misbehaving children, "the owls will get you", and in most Native American folklore, owls are a symbol of death. According to the Apache and Seminole tribes, hearing owls hooting is considered the subject of numerous "bogeyman" stories told to warn children to remain indoors at night or not to cry too much, otherwise the owl may carry them away. In some tribal legends, owls are associated with spirits of the dead, and the bony circles around an owl's eyes are said to comprise the fingernails of apparitional humans. Sometimes owls are said to carry messages from beyond the grave or deliver supernatural warnings to people who have broken tribal taboos. The Aztecs and the Maya, along with other natives of Mesoamerica, considered the owl a symbol of death and destruction. In fact, the Aztec god of death, Mictlantecuhtli, was often depicted with owls. There is an old saying in Mexico that is still in use: Cuando el tecolote canta, el indio muere ("When the owl cries/sings, the Indian dies"). The Popol Vuh, a Mayan religious text, describes owls as messengers of Xibalba (the Mayan "Place of Fright"). The belief that owls are messengers and harbingers of the dark powers is also found among the Hočągara (Winnebago) of Wisconsin. When in earlier days the Hočągara committed the sin of killing enemies while they were within the sanctuary of the chief's lodge, an owl appeared and spoke to them in the voice of a human, saying, "From now on, the Hočągara will have no luck." This marked the beginning of the decline of their tribe. An owl appeared to Glory of the Morning, the only female chief of the Hočąk nation, and uttered her name. Soon after, she died. According to the culture of the Hopi, a Uto-Aztec tribe, taboos surround owls, which are associated with sorcery and other evils. The Ojibwe tribes, as well as their Aboriginal Canadian counterparts, used an owl as a symbol for both evil and death. In addition, they used owls as a symbol of very high status of spiritual leaders of their spirituality. The Pawnee tribes viewed owls as the symbol of protection from any danger within their realms. The Puebloan peoples associated owls with Skeleton Man, the god of death and the spirit of fertility. The Yakama tribes use an owl as a totem, to guide where and how forests and natural resources are useful with management. Rodent control Encouraging natural predators to control rodent population is a natural form of pest control, along with excluding food sources for rodents. Placing a nest box for owls on a property can help control rodent populations (one family of hungry barn owls can consume more than 3,000 rodents in a nesting season) while maintaining the naturally balanced food chain. Attacks on humans Although humans and owls frequently live together in harmony, there have been incidents when owls have attacked humans. For example, in January 2013, a man from Inverness, Scotland suffered heavy bleeding and went into shock after being attacked by an owl, which was likely a eagle-owl. The photographer Eric Hosking lost his left eye after attempting to photograph a tawny owl, which inspired the title of his 1970 autobiography, An Eye for a Bird. Conservation issues Almost all owls are listed in Appendix II of the international CITES treaty (the Convention on Illegal Trade in Endangered Species of Wild Fauna and Flora) with four species listed in Appendix I. Although owls have long been hunted, a 2008 news story from Malaysia indicates that the magnitude of owl poaching may be on the rise. In November 2008, TRAFFIC reported the seizure of 900 plucked and "oven-ready" owls in Peninsular Malaysia. Said Chris Shepherd, Senior Programme Officer for TRAFFIC's Southeast Asia office, "This is the first time we know of where 'ready-prepared' owls have been seized in Malaysia, and it may mark the start of a new trend in wild meat from the region. We will be monitoring developments closely." TRAFFIC commended the Department of Wildlife and National Parks in Malaysia for the raid that exposed the huge haul of owls. Included in the seizure were dead and plucked barn owls, spotted wood owls, crested serpent eagles, barred eagles, and brown wood owls, as well as 7,000 live lizards. In addition to hunting, other threats to owl populations are habitat loss, pesticides, viruses, and vehicle collisions.	Biology and health sciences	Strigiformes	null
37674	https://en.wikipedia.org/wiki/Duck	Duck	Duck is the common name for numerous species of waterfowl in the family Anatidae. Ducks are generally smaller and shorter-necked than swans and geese, which are members of the same family. Divided among several subfamilies, they are a form taxon; they do not represent a monophyletic group (the group of all descendants of a single common ancestral species), since swans and geese are not considered ducks. Ducks are mostly aquatic birds, and may be found in both fresh water and sea water. Ducks are sometimes confused with several types of unrelated water birds with similar forms, such as loons or divers, grebes, gallinules and coots. Etymology The word duck comes from Old English 'diver', a derivative of the verb 'to duck, bend down low as if to get under something, or dive', because of the way many species in the dabbling duck group feed by upending; compare with Dutch and German 'to dive'. This word replaced Old English / 'duck', possibly to avoid confusion with other words, such as 'end' with similar forms. Other Germanic languages still have similar words for duck, for example, Dutch , German and Norwegian . The word / was inherited from Proto-Indo-European; cf. Latin anas "duck", Lithuanian 'duck', Ancient Greek / ( /) 'duck', and Sanskrit 'water bird', among others. A duckling is a young duck in downy plumage or baby duck, but in the food trade a young domestic duck which has just reached adult size and bulk and its meat is still fully tender, is sometimes labelled as a duckling. A male is called a drake and the female is called a duck, or in ornithology a hen. Taxonomy All ducks belong to the biological order Anseriformes, a group that contains the ducks, geese and swans, as well as the screamers, and the magpie goose. All except the screamers belong to the biological family Anatidae. Within the family, ducks are split into a variety of subfamilies and 'tribes'. The number and composition of these subfamilies and tribes is the cause of considerable disagreement among taxonomists. Some base their decisions on morphological characteristics, others on shared behaviours or genetic studies. The number of suggested subfamilies containing ducks ranges from two to five. The significant level of hybridisation that occurs among wild ducks complicates efforts to tease apart the relationships between various species. In most modern classifications, the so-called 'true ducks' belong to the subfamily Anatinae, which is further split into a varying number of tribes. The largest of these, the Anatini, contains the 'dabbling' or 'river' ducks – named for their method of feeding primarily at the surface of fresh water. The 'diving ducks', also named for their primary feeding method, make up the tribe Aythyini. The 'sea ducks' of the tribe Mergini are diving ducks which specialise on fish and shellfish and spend a majority of their lives in saltwater. The tribe Oxyurini contains the 'stifftails', diving ducks notable for their small size and stiff, upright tails. A number of other species called ducks are not considered to be 'true ducks', and are typically placed in other subfamilies or tribes. The whistling ducks are assigned either to a tribe (Dendrocygnini) in the subfamily Anatinae or the subfamily Anserinae, or to their own subfamily (Dendrocygninae) or family (Dendrocyganidae). The freckled duck of Australia is either the sole member of the tribe Stictonettini in the subfamily Anserinae, or in its own family, the Stictonettinae. The shelducks make up the tribe Tadornini in the family Anserinae in some classifications, and their own subfamily, Tadorninae, in others, while the steamer ducks are either placed in the family Anserinae in the tribe Tachyerini or lumped with the shelducks in the tribe Tadorini. The perching ducks make up in the tribe Cairinini in the subfamily Anserinae in some classifications, while that tribe is eliminated in other classifications and its members assigned to the tribe Anatini. The torrent duck is generally included in the subfamily Anserinae in the monotypic tribe Merganettini, but is sometimes included in the tribe Tadornini. The pink-eared duck is sometimes included as a true duck either in the tribe Anatini or the tribe Malacorhynchini, and other times is included with the shelducks in the tribe Tadornini. Morphology The overall body plan of ducks is elongated and broad, and they are also relatively long-necked, albeit not as long-necked as the geese and swans. The body shape of diving ducks varies somewhat from this in being more rounded. The bill is usually broad and contains serrated pectens, which are particularly well defined in the filter-feeding species. In the case of some fishing species the bill is long and strongly serrated. The scaled legs are strong and well developed, and generally set far back on the body, more so in the highly aquatic species. The wings are very strong and are generally short and pointed, and the flight of ducks requires fast continuous strokes, requiring in turn strong wing muscles. Three species of steamer duck are almost flightless, however. Many species of duck are temporarily flightless while moulting; they seek out protected habitat with good food supplies during this period. This moult typically precedes migration. The drakes of northern species often have extravagant plumage, but that is moulted in summer to give a more female-like appearance, the "eclipse" plumage. Southern resident species typically show less sexual dimorphism, although there are exceptions such as the paradise shelduck of New Zealand, which is both strikingly sexually dimorphic and in which the female's plumage is brighter than that of the male. The plumage of juvenile birds generally resembles that of the female. Female ducks have evolved to have a corkscrew shaped vagina to prevent rape. Distribution and habitat Ducks have a cosmopolitan distribution, and are found on every continent except Antarctica. Several species manage to live on subantarctic islands, including South Georgia and the Auckland Islands. Ducks have reached a number of isolated oceanic islands, including the Hawaiian Islands, Micronesia and the Galápagos Islands, where they are often and less often . A handful are endemic to such far-flung islands. Some duck species, mainly those breeding in the temperate and Arctic Northern Hemisphere, are migratory; those in the tropics are generally not. Some ducks, particularly in Australia where rainfall is erratic, are nomadic, seeking out the temporary lakes and pools that form after localised heavy rain. Behaviour Feeding Ducks eat food sources such as grasses, aquatic plants, fish, insects, small amphibians, worms, and small molluscs. Dabbling ducks feed on the surface of water or on land, or as deep as they can reach by up-ending without completely submerging. Along the edge of the bill, there is a comb-like structure called a pecten. This strains the water squirting from the side of the bill and traps any food. The pecten is also used to preen feathers and to hold slippery food items. Diving ducks and sea ducks forage deep underwater. To be able to submerge more easily, the diving ducks are heavier than dabbling ducks, and therefore have more difficulty taking off to fly. A few specialized species such as the mergansers are adapted to catch and swallow large fish. The others have the characteristic wide flat bill adapted to dredging-type jobs such as pulling up waterweed, pulling worms and small molluscs out of mud, searching for insect larvae, and bulk jobs such as dredging out, holding, turning head first, and swallowing a squirming frog. To avoid injury when digging into sediment it has no cere, but the nostrils come out through hard horn. The Guardian published an article advising that ducks should not be fed with bread because it damages the health of the ducks and pollutes waterways. Breeding Ducks generally only have one partner at a time, although the partnership usually only lasts one year. Larger species and the more sedentary species (like fast-river specialists) tend to have pair-bonds that last numerous years. Most duck species breed once a year, choosing to do so in favourable conditions (spring/summer or wet seasons). Ducks also tend to make a nest before breeding, and, after hatching, lead their ducklings to water. Mother ducks are very caring and protective of their young, but may abandon some of their ducklings if they are physically stuck in an area they cannot get out of (such as nesting in an enclosed courtyard) or are not prospering due to genetic defects or sickness brought about by hypothermia, starvation, or disease. Ducklings can also be orphaned by inconsistent late hatching where a few eggs hatch after the mother has abandoned the nest and led her ducklings to water. Communication Female mallard ducks (as well as several other species in the genus Anas, such as the American and Pacific black ducks, spot-billed duck, northern pintail and common teal) make the classic "quack" sound while males make a similar but raspier sound that is sometimes written as "breeeeze", but, despite widespread misconceptions, most species of duck do not "quack". In general, ducks make a range of calls, including whistles, cooing, yodels and grunts. For example, the scaup – which are diving ducks – make a noise like "scaup" (hence their name). Calls may be loud displaying calls or quieter contact calls. A common urban legend claims that duck quacks do not echo; however, this has been proven to be false. This myth was first debunked by the Acoustics Research Centre at the University of Salford in 2003 as part of the British Association's Festival of Science. It was also debunked in one of the earlier episodes of the popular Discovery Channel television show MythBusters. Predators Ducks have many predators. Ducklings are particularly vulnerable, since their inability to fly makes them easy prey not only for predatory birds but also for large fish like pike, crocodilians, predatory testudines such as the alligator snapping turtle, and other aquatic hunters, including fish-eating birds such as herons. Ducks' nests are raided by land-based predators, and brooding females may be caught unaware on the nest by mammals, such as foxes, or large birds, such as hawks or owls. Adult ducks are fast fliers, but may be caught on the water by large aquatic predators including big fish such as the North American muskie and the European pike. In flight, ducks are safe from all but a few predators such as humans and the peregrine falcon, which uses its speed and strength to catch ducks. Relationship with humans Hunting Humans have hunted ducks since prehistoric times. Excavations of middens in California dating to 7800 – 6400 BP have turned up bones of ducks, including at least one now-extinct flightless species. Ducks were captured in "significant numbers" by Holocene inhabitants of the lower Ohio River valley, suggesting they took advantage of the seasonal bounty provided by migrating waterfowl. Neolithic hunters in locations as far apart as the Caribbean, Scandinavia, Egypt, Switzerland, and China relied on ducks as a source of protein for some or all of the year. Archeological evidence shows that Māori people in New Zealand hunted the flightless Finsch's duck, possibly to extinction, though rat predation may also have contributed to its fate. A similar end awaited the Chatham duck, a species with reduced flying capabilities which went extinct shortly after its island was colonised by Polynesian settlers. It is probable that duck eggs were gathered by Neolithic hunter-gathers as well, though hard evidence of this is uncommon. In many areas, wild ducks (including ducks farmed and released into the wild) are hunted for food or sport, by shooting, or by being trapped using duck decoys. Because an idle floating duck or a duck squatting on land cannot react to fly or move quickly, "a sitting duck" has come to mean "an easy target". These ducks may be contaminated by pollutants such as PCBs. Domestication Ducks have many economic uses, being farmed for their meat, eggs, and feathers (particularly their down). Approximately 3 billion ducks are slaughtered each year for meat worldwide. They are also kept and bred by aviculturists and often displayed in zoos. Almost all the varieties of domestic ducks are descended from the mallard (Anas platyrhynchos), apart from the Muscovy duck (Cairina moschata). The Call duck is another example of a domestic duck breed. Its name comes from its original use established by hunters, as a decoy to attract wild mallards from the sky, into traps set for them on the ground. The call duck is the world's smallest domestic duck breed, as it weighs less than . Heraldry Ducks appear on several coats of arms, including the coat of arms of Lubāna (Latvia) and the coat of arms of Föglö (Åland). Cultural references In 2002, psychologist Richard Wiseman and colleagues at the University of Hertfordshire, UK, finished a year-long LaughLab experiment, concluding that of all animals, ducks attract the most humor and silliness; he said, "If you're going to tell a joke involving an animal, make it a duck." The word "duck" may have become an inherently funny word in many languages, possibly because ducks are seen as silly in their looks or behavior. Of the many ducks in fiction, many are cartoon characters, such as Walt Disney's Donald Duck, and Warner Bros.' Daffy Duck. Howard the Duck started as a comic book character in 1973 and was made into a movie in 1986. The 1992 Disney film The Mighty Ducks, starring Emilio Estevez, chose the duck as the mascot for the fictional youth hockey team who are protagonists of the movie, based on the duck being described as a fierce fighter. This led to the duck becoming the nickname and mascot for the eventual National Hockey League professional team of the Anaheim Ducks, who were founded with the name the Mighty Ducks of Anaheim. The duck is also the nickname of the University of Oregon sports teams as well as the Long Island Ducks minor league baseball team.	Biology and health sciences	Anseriformes	null
37694	https://en.wikipedia.org/wiki/Seed	Seed	In botany, a seed is a plant embryo and nutrient reserve enclosed in a seed coat, a protective outer covering called a testa. More generally, the term "seed" means anything that can be sown, which may include seed and husk or tuber. Seeds are the product of the ripened ovule, after the embryo sac is fertilized by sperm from pollen, forming a zygote. The embryo within a seed develops from the zygote and grows within the mother plant to a certain size before growth is halted. The formation of the seed is the defining part of the process of reproduction in seed plants (spermatophytes). Other plants such as ferns, mosses and liverworts, do not have seeds and use water-dependent means to propagate themselves. Seed plants now dominate biological niches on land, from forests to grasslands both in hot and cold climates. In the flowering plants, the ovary ripens into a fruit which contains the seed and serves to disseminate it. Many structures commonly referred to as "seeds" are actually dry fruits. Sunflower seeds are sometimes sold commercially while still enclosed within the hard wall of the fruit, which must be split open to reach the seed. Different groups of plants have other modifications, the so-called stone fruits (such as the peach) have a hardened fruit layer (the endocarp) fused to and surrounding the actual seed. Nuts are the one-seeded, hard-shelled fruit of some plants with an indehiscent seed, such as an acorn or hazelnut. History The first land plants evolved around 468 million years ago, and reproduced using spores. The earliest seed bearing plants to appear were the gymnosperms, which have no ovaries to contain the seeds. They arose during the late Devonian period (416 million to 358 million years ago). From these early gymnosperms, seed ferns evolved during the Carboniferous period (359 to 299 million years ago); they had ovules that were borne in a cupule, which consisted of groups of enclosing branches likely used to protect the developing seed. Published literature about seed storage, viability and its hygrometric dependence began in the early 19th century, influential works being: 1832 seed storage guide in Augustin Pyramus de Candolle's Conservation des Graines, part of his 3-volume Physiologie végétale, ou Exposition des forces et des fonctions vitales des végétaux (1832, v. 2, pp. 618–626, Paris); (translated title, "Plant physiology, or Exposition of the vital forces and functions of plants") 1846 viability studies by Augustin de Candolle, published in "Sur la durée relative de la faculté de germer des graines appartenant à diverses familles" (Annales des Sciences Naturelles; Botanique, 1846, III 6: 373–382); (translated title, "On the relative duration of the ability to germinate seeds belonging to various families") 1897 seed hygrometric studies by Victor Jodin (Annales Agronomiques, October 1897) 1912's Henry B. Guppy's 528 page "Studies in Seeds and Fruits- An Investigation with the Balance" (1912, London, England); subsequently reviewed in Science (June 1914, Washington, D.C.) Development Angiosperm seeds are "enclosed seeds", produced in a hard or fleshy structure called a fruit that encloses them for protection. Some fruits have layers of both hard and fleshy material. In gymnosperms, no special structure develops to enclose the seeds, which begin their development "naked" on the bracts of cones. However, the seeds do become covered by the cone scales as they develop in some species of conifer. Angiosperm (flowering plants) seeds consist of three genetically distinct constituents: (1) the embryo formed from the zygote, (2) the endosperm, which is normally triploid, (3) the seed coat from tissue derived from the maternal tissue of the ovule. In angiosperms, the process of seed development begins with double fertilization, which involves the fusion of two male gametes with the egg cell and the central cell to form the primary endosperm and the zygote. Right after fertilization, the zygote is mostly inactive, but the primary endosperm divides rapidly to form the endosperm tissue. This tissue becomes the food the young plant will consume until the roots have developed after germination. Ovule After fertilization, the ovules develop into the seeds. The ovule consists of a number of components: The funicle (funiculus, funiculi) or seed stalk which attaches the ovule to the placenta and hence ovary or fruit wall, at the pericarp. The nucellus, the remnant of the megasporangium and main region of the ovule where the megagametophyte develops. The micropyle, a small pore or opening in the apex of the integument of the ovule where the pollen tube usually enters during the process of fertilization. The chalaza, the base of the ovule opposite the micropyle, where integument and nucellus are joined. The shape of the ovules as they develop often affects the final shape of the seeds. Plants generally produce ovules of four shapes: the most common shape is called anatropous, with a curved shape. Orthotropous ovules are straight with all the parts of the ovule lined up in a long row producing an uncurved seed. Campylotropous ovules have a curved megagametophyte often giving the seed a tight "C" shape. The last ovule shape is called amphitropous, where the ovule is partly inverted and turned back 90 degrees on its stalk (the funicle or funiculus). In the majority of flowering plants, the zygote's first division is transversely oriented in regards to the long axis, and this establishes the polarity of the embryo. The upper or chalazal pole becomes the main area of growth of the embryo, while the lower or micropylar pole produces the stalk-like suspensor that attaches to the micropyle. The suspensor absorbs and manufactures nutrients from the endosperm that are used during the embryo's growth. Embryo The main components of the embryo are: The cotyledons, the seed leaves, attached to the embryonic axis. There may be one (Monocotyledons), or two (Dicotyledons). The cotyledons are also the source of nutrients in the non-endospermic dicotyledons, in which case they replace the endosperm, and are thick and leathery. In endospermic seeds, the cotyledons are thin and papery. Dicotyledons have the point of attachment opposite one another on the axis. The epicotyl, the embryonic axis above the point of attachment of the cotyledon(s). The plumule, the tip of the epicotyl, and has a feathery appearance due to the presence of young leaf primordia at the apex, and will become the shoot upon germination. The hypocotyl, the embryonic axis below the point of attachment of the cotyledon(s), connecting the epicotyl and the radicle, being the stem-root transition zone. The radicle, the basal tip of the hypocotyl, grows into the primary root. Monocotyledonous plants have two additional structures in the form of sheaths. The plumule is covered with a coleoptile that forms the first leaf while the radicle is covered with a coleorhiza that connects to the primary root and adventitious roots form the sides. Here the hypocotyl is a rudimentary axis between radicle and plumule. The seeds of corn are constructed with these structures; pericarp, scutellum (single large cotyledon) that absorbs nutrients from the endosperm, plumule, radicle, coleoptile, and coleorhiza – these last two structures are sheath-like and enclose the plumule and radicle, acting as a protective covering. Seed coat The maturing ovule undergoes marked changes in the integuments, generally a reduction and disorganization but occasionally a thickening. The seed coat forms from the two integuments or outer layers of cells of the ovule, which derive from tissue from the mother plant, the inner integument forms the tegmen and the outer forms the testa. (The seed coats of some monocotyledon plants, such as the grasses, are not distinct structures, but are fused with the fruit wall to form a pericarp.) The testae of both monocots and dicots are often marked with patterns and textured markings, or have wings or tufts of hair. When the seed coat forms from only one layer, it is also called the testa, though not all such testae are homologous from one species to the next. The funiculus abscisses (detaches at fixed point – abscission zone), the scar forming an oval depression, the hilum. Anatropous ovules have a portion of the funiculus that is adnate (fused to the seed coat), and which forms a longitudinal ridge, or raphe, just above the hilum. In bitegmic ovules (e.g. Gossypium described here) both inner and outer integuments contribute to the seed coat formation. With continuing maturation the cells enlarge in the outer integument. While the inner epidermis may remain a single layer, it may also divide to produce two to three layers and accumulates starch, and is referred to as the colourless layer. By contrast, the outer epidermis becomes tanniferous. The inner integument may consist of eight to fifteen layers. As the cells enlarge, and starch is deposited in the outer layers of the pigmented zone below the outer epidermis, this zone begins to lignify, while the cells of the outer epidermis enlarge radially and their walls thicken, with nucleus and cytoplasm compressed into the outer layer. these cells which are broader on their inner surface are called palisade cells. In the inner epidermis, the cells also enlarge radially with plate like thickening of the walls. The mature inner integument has a palisade layer, a pigmented zone with 15–20 layers, while the innermost layer is known as the fringe layer. Gymnosperms In gymnosperms, which do not form ovaries, the ovules and hence the seeds are exposed. This is the basis for their nomenclature – naked seeded plants. Two sperm cells transferred from the pollen do not develop the seed by double fertilization, but one sperm nucleus unites with the egg nucleus and the other sperm is not used. Sometimes each sperm fertilizes an egg cell and one zygote is then aborted or absorbed during early development. The seed is composed of the embryo (the result of fertilization) and tissue from the mother plant, which also form a cone around the seed in coniferous plants such as pine and spruce. Shape and appearance Seeds are very diverse, and as such there are many terms are used to describe them. Terms to describe shape Bean-shaped () – resembling a kidney, with lobed ends on either side of the hilum Square or Oblong – angular, with all sides being either equal, or longer-than-wide Triangular – three-sided, broadest below the middle Elliptic or Ovate or Obovate – rounded at both ends, or egg shaped (ovate or obovate, broader at one end), being rounded but either symmetrical about the middle, or broader below the middle, or broader above the middle Discoid – resembling a disc or plate, having both thickness and parallel faces and with a rounded margin) Ellipsoid Globose – spherical Subglobose (Inflated, but less than spherical) Lenticular Ovoid Sectoroid Other common descriptors for seeds focus on color, texture, and form. Striate seeds are striped with parallel, longitudinal lines or ridges. The most common colours are brown and black, with other colours appearing less frequently. The surface texture varies from highly polished to considerably roughened. The surface may also have a variety of appendages (see Seed coat), and be described by terms such as papillate or digitiform (finger-like). A seed coat with the consistency of cork is referred to as suberose. Other terms include crustaceous (hard, thin or brittle). Structure A typical seed includes two basic parts: an embryo; a seed coat. In addition, the endosperm forms a supply of nutrients for the embryo in most monocotyledons and the endospermic dicotyledons. Seed types Seeds have been considered to occur in many structurally different types (Martin 1946). These are based on a number of criteria, of which the dominant one is the embryo-to-seed size ratio. This reflects the degree to which the developing cotyledons absorb the nutrients of the endosperm, and thus obliterate it. Six types occur amongst the monocotyledons, ten in the dicotyledons, and two in the gymnosperms (linear and spatulate). This classification is based on three characteristics: embryo morphology, amount of endosperm and the position of the embryo relative to the endosperm. Embryo In endospermic seeds, there are two distinct regions inside the seed coat, an upper and larger endosperm and a lower smaller embryo. The embryo is the fertilised ovule, an immature plant from which a new plant will grow under proper conditions. The embryo has one cotyledon or seed leaf in monocotyledons, two cotyledons in almost all dicotyledons and two or more in gymnosperms. In the fruit of grains (caryopses) the single monocotyledon is shield shaped and hence called a scutellum. The scutellum is pressed closely against the endosperm from which it absorbs food and passes it to the growing parts. Embryo descriptors include small, straight, bent, curved, and curled. Nutrient storage Within the seed, there usually is a store of nutrients for the seedling that will grow from the embryo. The form of the stored nutrition varies depending on the kind of plant. In angiosperms, the stored food begins as a tissue called the endosperm, which is derived from the mother plant and the pollen via double fertilization. It is usually triploid, and is rich in oil or starch, and protein. In gymnosperms, such as conifers, the food storage tissue (also called endosperm) is part of the female gametophyte, a haploid tissue. The endosperm is surrounded by the aleurone layer (peripheral endosperm), filled with proteinaceous aleurone grains. Originally, by analogy with the animal ovum, the outer nucellus layer (perisperm) was referred to as albumen, and the inner endosperm layer as vitellus. Although misleading, the term began to be applied to all the nutrient matter. This terminology persists in referring to endospermic seeds as "albuminous". The nature of this material is used in both describing and classifying seeds, in addition to the embryo to endosperm size ratio. The endosperm may be considered to be farinaceous (or mealy) in which the cells are filled with starch, as for instance cereal grains, or not (non-farinaceous). The endosperm may also be referred to as "fleshy" or "cartilaginous" with thicker soft cells such as coconut, but may also be oily as in Ricinus (castor oil), Croton and Poppy. The endosperm is called "horny" when the cell walls are thicker such as date and coffee, or "ruminated" if mottled, as in nutmeg, palms and Annonaceae. In most monocotyledons (such as grasses and palms) and some (endospermic or albuminous) dicotyledons (such as castor beans) the embryo is embedded in the endosperm (and nucellus), which the seedling will use upon germination. In the non-endospermic dicotyledons the endosperm is absorbed by the embryo as the latter grows within the developing seed, and the cotyledons of the embryo become filled with stored food. At maturity, seeds of these species have no endosperm and are also referred to as exalbuminous seeds. The exalbuminous seeds include the legumes (such as beans and peas), trees such as the oak and walnut, vegetables such as squash and radish, and sunflowers. According to Bewley and Black (1978), Brazil nut storage is in hypocotyl and this place of storage is uncommon among seeds. All gymnosperm seeds are albuminous. Seed coat The seed coat develops from the maternal tissue, the integuments, originally surrounding the ovule. The seed coat in the mature seed can be a paper-thin layer (e.g. peanut) or something more substantial (e.g. thick and hard in honey locust and coconut), or fleshy as in the sarcotesta of pomegranate. The seed coat helps protect the embryo from mechanical injury, predators, and drying out. Depending on its development, the seed coat is either bitegmic or unitegmic. Bitegmic seeds form a testa from the outer integument and a tegmen from the inner integument while unitegmic seeds have only one integument. Usually, parts of the testa or tegmen form a hard protective mechanical layer. The mechanical layer may prevent water penetration and germination. Amongst the barriers may be the presence of lignified sclereids. The outer integument has a number of layers, generally between four and eight organised into three layers: (a) outer epidermis, (b) outer pigmented zone of two to five layers containing tannin and starch, and (c) inner epidermis. The endotegmen is derived from the inner epidermis of the inner integument, the exotegmen from the outer surface of the inner integument. The endotesta is derived from the inner epidermis of the outer integument, and the outer layer of the testa from the outer surface of the outer integument is referred to as the exotesta. If the exotesta is also the mechanical layer, this is called an exotestal seed, but if the mechanical layer is the endotegmen, then the seed is endotestal. The exotesta may consist of one or more rows of cells that are elongated and pallisade like (e.g. Fabaceae), hence 'palisade exotesta'. In addition to the three basic seed parts, some seeds have an appendage, an aril, a fleshy outgrowth of the funicle (funiculus), (as in yew and nutmeg) or an oily appendage, an elaiosome (as in Corydalis), or hairs (trichomes). In the latter example these hairs are the source of the textile crop cotton. Other seed appendages include the raphe (a ridge), wings, caruncles (a soft spongy outgrowth from the outer integument in the vicinity of the micropyle), spines, or tubercles. A scar also may remain on the seed coat, called the hilum, where the seed was attached to the ovary wall by the funicle. Just below it is a small pore, representing the micropyle of the ovule. Size and seed set Seeds are very diverse in size. The dust-like orchid seeds are the smallest, with about one million seeds per gram; they are often embryonic seeds with immature embryos and no significant energy reserves. Orchids and a few other groups of plants are mycoheterotrophs which depend on mycorrhizal fungi for nutrition during germination and the early growth of the seedling. Some terrestrial orchid seedlings, in fact, spend the first few years of their lives deriving energy from the fungi and do not produce green leaves. At up to 55 pounds (25 kilograms) the largest seed is the coco de mer(Lodoicea maldivica). This indicates a 25 Billion fold difference in seed weight. Plants that produce smaller seeds can generate many more seeds per flower, while plants with larger seeds invest more resources into those seeds and normally produce fewer seeds. Small seeds are quicker to ripen and can be dispersed sooner, so autumn all blooming plants often have small seeds. Many annual plants produce great quantities of smaller seeds; this helps to ensure at least a few will end in a favorable place for growth. Herbaceous perennials and woody plants often have larger seeds; they can produce seeds over many years, and larger seeds have more energy reserves for germination and seedling growth and produce larger, more established seedlings after germination. Functions Seeds serve several functions for the plants that produce them. Key among these functions are nourishment of the embryo, dispersal to a new location, and dormancy during unfavorable conditions. Seeds fundamentally are means of reproduction, and most seeds are the product of sexual reproduction which produces a remixing of genetic material and phenotype variability on which natural selection acts. Plant seeds hold endophytic microorganisms that can perform various functions, the most important of which is protection against disease. Embryo nourishment Seeds protect and nourish the embryo or young plant. They usually give a seedling a faster start than a sporeling from a spore, because of the larger food reserves in the seed and the multicellularity of the enclosed embryo. Dispersal Unlike animals, plants are limited in their ability to seek out favorable conditions for life and growth. As a result, plants have evolved many ways to disperse their offspring by dispersing their seeds (see also vegetative reproduction). A seed must somehow "arrive" at a location and be there at a time favorable for germination and growth. When the fruits open and release their seeds in a regular way, it is called dehiscent, which is often distinctive for related groups of plants; these fruits include capsules, follicles, legumes, silicles and siliques. When fruits do not open and release their seeds in a regular fashion, they are called indehiscent, which include the fruits achenes, caryopses, nuts, samaras, and utricles. By wind (anemochory) Some seeds (e.g., pine) have a wing that aids in wind dispersal. The dustlike seeds of orchids are carried efficiently by the wind. Some seeds (e.g. milkweed, poplar) have hairs that aid in wind dispersal. Other seeds are enclosed in fruit structures that aid wind dispersal in similar ways: Dandelion achenes have hairs. Maple samaras have two wings. By water (hydrochory) Some plants, such as Mucuna and Dioclea, produce buoyant seeds termed sea-beans or drift seeds because they float in rivers to the oceans and wash up on beaches. By animals (zoochory) Seeds (burrs) with barbs or hooks (e.g. acaena, burdock, dock) which attach to animal fur or feathers, and then drop off later. Seeds with a fleshy covering (e.g. apple, cherry, juniper) are eaten by animals (birds, mammals, reptiles, fish) which then disperse these seeds in their droppings. Seeds (nuts) are attractive long-term storable food resources for animals (e.g. acorns, hazelnut, walnut); the seeds are stored some distance from the parent plant, and some escape being eaten if the animal forgets them. Myrmecochory is the dispersal of seeds by ants. Foraging ants disperse seeds which have appendages called elaiosomes (e.g. bloodroot, trilliums, acacias, and many species of Proteaceae). Elaiosomes are soft, fleshy structures that contain nutrients for animals that eat them. The ants carry such seeds back to their nest, where the elaiosomes are eaten. The remainder of the seed, which is hard and inedible to the ants, then germinates either within the nest or at a removal site where the seed has been discarded by the ants. This dispersal relationship is an example of mutualism, since the plants depend upon the ants to disperse seeds, while the ants depend upon the plants seeds for food. As a result, a drop in numbers of one partner can reduce success of the other. In South Africa, the Argentine ant (Linepithema humile) has invaded and displaced native species of ants. Unlike the native ant species, Argentine ants do not collect the seeds of Mimetes cucullatus or eat the elaiosomes. In areas where these ants have invaded, the numbers of Mimetes seedlings have dropped. Dormancy Seed dormancy has two main functions: the first is synchronizing germination with the optimal conditions for survival of the resulting seedling; the second is spreading germination of a batch of seeds over time so a catastrophe (e.g. late frosts, drought, herbivory) does not result in the death of all offspring of a plant (bet-hedging). Seed dormancy is defined as a seed failing to germinate under environmental conditions optimal for germination, normally when the environment is at a suitable temperature with proper soil moisture. This true dormancy or innate dormancy is therefore caused by conditions within the seed that prevent germination. Thus dormancy is a state of the seed, not of the environment. Induced dormancy, enforced dormancy or seed quiescence occurs when a seed fails to germinate because the external environmental conditions are inappropriate for germination, mostly in response to conditions being too dark or light, too cold or hot, or too dry. Seed dormancy is not the same as seed persistence in the soil or on the plant, though even in scientific publications dormancy and persistence are often confused or used as synonyms. Often, seed dormancy is divided into four major categories: exogenous; endogenous; combinational; and secondary. A more recent system distinguishes five classes: morphological, physiological, morphophysiological, physical, and combinational dormancy. Exogenous dormancy is caused by conditions outside the embryo, including: Physical dormancy or hard seed coats occurs when seeds are impermeable to water. At dormancy break, a specialized structure, the 'water gap', is disrupted in response to environmental cues, especially temperature, so water can enter the seed and germination can occur. Plant families where physical dormancy occurs include Anacardiaceae, Cannaceae, Convulvulaceae, Fabaceae and Malvaceae. Chemical dormancy considers species that lack physiological dormancy, but where a chemical prevents germination. This chemical can be leached out of the seed by rainwater or snow melt or be deactivated somehow. Leaching of chemical inhibitors from the seed by rain water is often cited as an important cause of dormancy release in seeds of desert plants, but little evidence exists to support this claim. Endogenous dormancy is caused by conditions within the embryo itself, including: In morphological dormancy, germination is prevented due to morphological characteristics of the embryo. In some species, the embryo is just a mass of cells when seeds are dispersed; it is not differentiated. Before germination can take place, both differentiation and growth of the embryo have to occur. In other species, the embryo is differentiated but not fully grown (underdeveloped) at dispersal, and embryo growth up to a species specific length is required before germination can occur. Examples of plant families where morphological dormancy occurs are Apiaceae, Cycadaceae, Liliaceae, Magnoliaceae and Ranunculaceae. Morphophysiological dormancy includes seeds with underdeveloped embryos, and also have physiological components to dormancy. These seeds, therefore, require a dormancy-breaking treatments, as well as a period of time to develop fully grown embryos. Plant families where morphophysiological dormancy occurs include Apiaceae, Aquifoliaceae, Liliaceae, Magnoliaceae, Papaveraceae and Ranunculaceae. Some plants with morphophysiological dormancy, such as Asarum or Trillium species, have multiple types of dormancy, one affects radicle (root) growth, while the other affects plumule (shoot) growth. The terms "double dormancy" and "two-year seeds" are used for species whose seeds need two years to complete germination or at least two winters and one summer. Dormancy of the radicle (seedling root) is broken during the first winter after dispersal while dormancy of the shoot bud is broken during the second winter. Physiological dormancy means the embryo, due to physiological causes, cannot generate enough power to break through the seed coat, endosperm or other covering structures. Dormancy is typically broken at cool wet, warm wet, or warm dry conditions. Abscisic acid is usually the growth inhibitor in seeds, and its production can be affected by light. Drying, in some plants, including a number of grasses and those from seasonally arid regions, is needed before they will germinate. The seeds are released, but need to have a lower moisture content before germination can begin. If the seeds remain moist after dispersal, germination can be delayed for many months or even years. Many herbaceous plants from temperate climate zones have physiological dormancy that disappears with drying of the seeds. Other species will germinate after dispersal only under very narrow temperature ranges, but as the seeds dry, they are able to germinate over a wider temperature range. In seeds with combinational dormancy, the seed or fruit coat is impermeable to water and the embryo has physiological dormancy. Depending on the species, physical dormancy can be broken before or after physiological dormancy is broken. Secondary dormancy* is caused by conditions after the seed has been dispersed and occurs in some seeds when nondormant seed is exposed to conditions that are not favorable to germination, very often high temperatures. The mechanisms of secondary dormancy are not yet fully understood, but might involve the loss of sensitivity in receptors in the plasma membrane. The following types of seed dormancy do not involve seed dormancy, strictly speaking, as lack of germination is prevented by the environment, not by characteristics of the seed itself (see Germination): Photodormancy or light sensitivity affects germination of some seeds. These photoblastic seeds need a period of darkness or light to germinate. In species with thin seed coats, light may be able to penetrate into the dormant embryo. The presence of light or the absence of light may trigger the germination process, inhibiting germination in some seeds buried too deeply or in others not buried in the soil. Thermodormancy is seed sensitivity to heat or cold. Some seeds, including cocklebur and amaranth, germinate only at high temperatures (30 °C or 86 °F); many plants that have seeds that germinate in early to midsummer have thermodormancy, so germinate only when the soil temperature is warm. Other seeds need cool soils to germinate, while others, such as celery, are inhibited when soil temperatures are too warm. Often, thermodormancy requirements disappear as the seed ages or dries. Not all seeds undergo a period of dormancy. Seeds of some mangroves are viviparous; they begin to germinate while still attached to the parent. The large, heavy root allows the seed to penetrate into the ground when it falls. Many garden plant seeds will germinate readily as soon as they have water and are warm enough; though their wild ancestors may have had dormancy, these cultivated plants lack it. After many generations of selective pressure by plant breeders and gardeners, dormancy has been selected out. For annuals, seeds are a way for the species to survive dry or cold seasons. Ephemeral plants are usually annuals that can go from seed to seed in as few as six weeks. Persistence and seed banks Germination Seed germination is a process by which a seed embryo develops into a seedling. It involves the reactivation of the metabolic pathways that lead to growth and the emergence of the radicle or seed root and plumule or shoot. The emergence of the seedling above the soil surface is the next phase of the plant's growth and is called seedling establishment. Three fundamental conditions must exist before germination can occur. (1) The embryo must be alive, called seed viability. (2) Any dormancy requirements that prevent germination must be overcome. (3) The proper environmental conditions must exist for germination. Far red light can prevent germination. Seed viability is the ability of the embryo to germinate and is affected by a number of different conditions. Some plants do not produce seeds that have functional complete embryos, or the seed may have no embryo at all, often called empty seeds. Predators and pathogens can damage or kill the seed while it is still in the fruit or after it is dispersed. Environmental conditions like flooding or heat can kill the seed before or during germination. The age of the seed affects its health and germination ability: since the seed has a living embryo, over time cells die and cannot be replaced. Some seeds can live for a long time before germination, while others can only survive for a short period after dispersal before they die. Seed vigor is a measure of the quality of seed, and involves the viability of the seed, the germination percentage, germination rate, and the strength of the seedlings produced. The germination percentage is simply the proportion of seeds that germinate from all seeds subject to the right conditions for growth. The germination rate is the length of time it takes for the seeds to germinate. Germination percentages and rates are affected by seed viability, dormancy and environmental effects that impact on the seed and seedling. In agriculture and horticulture quality seeds have high viability, measured by germination percentage plus the rate of germination. This is given as a percent of germination over a certain amount of time, 90% germination in 20 days, for example. 'Dormancy' is covered above; many plants produce seeds with varying degrees of dormancy, and different seeds from the same fruit can have different degrees of dormancy. It's possible to have seeds with no dormancy if they are dispersed right away and do not dry (if the seeds dry they go into physiological dormancy). There is great variation amongst plants and a dormant seed is still a viable seed even though the germination rate might be very low. Environmental conditions affecting seed germination include; water, oxygen, temperature and light. Three distinct phases of seed germination occur: water imbibition; lag phase; and radicle emergence. In order for the seed coat to split, the embryo must imbibe (soak up water), which causes it to swell, splitting the seed coat. However, the nature of the seed coat determines how rapidly water can penetrate and subsequently initiate germination. The rate of imbibition is dependent on the permeability of the seed coat, amount of water in the environment and the area of contact the seed has to the source of water. For some seeds, imbibing too much water too quickly can kill the seed. For some seeds, once water is imbibed the germination process cannot be stopped, and drying then becomes fatal. Other seeds can imbibe and lose water a few times without causing ill effects, but drying can cause secondary dormancy. Repair of DNA damage During seed dormancy, often associated with unpredictable and stressful environments, DNA damage accumulates as the seeds age. In rye seeds, the reduction of DNA integrity due to damage is associated with loss of seed viability during storage. Upon germination, seeds of Vicia faba undergo DNA repair. A plant DNA ligase that is involved in repair of single- and double-strand breaks during seed germination is an important determinant of seed longevity. Also, in Arabidopsis seeds, the activities of the DNA repair enzymes Poly ADP ribose polymerases (PARP) are likely needed for successful germination. Thus DNA damages that accumulate during dormancy appear to be a problem for seed survival, and the enzymatic repair of DNA damages during germination appears to be important for seed viability. Inducing germination A number of different strategies are used by gardeners and horticulturists to break seed dormancy. Scarification allows water and gases to penetrate into the seed; it includes methods to physically break the hard seed coats or soften them by chemicals, such as soaking in hot water or poking holes in the seed with a pin or rubbing them on sandpaper or cracking with a press or hammer. Sometimes fruits are harvested while the seeds are still immature and the seed coat is not fully developed and sown right away before the seed coat become impermeable. Under natural conditions, seed coats are worn down by rodents chewing on the seed, the seeds rubbing against rocks (seeds are moved by the wind or water currents), by undergoing freezing and thawing of surface water, or passing through an animal's digestive tract. In the latter case, the seed coat protects the seed from digestion, while often weakening the seed coat such that the embryo is ready to sprout when it is deposited, along with a bit of fecal matter that acts as fertilizer, far from the parent plant. Microorganisms are often effective in breaking down hard seed coats and are sometimes used by people as a treatment; the seeds are stored in a moist warm sandy medium for several months under nonsterile conditions. Stratification, also called moist-chilling, breaks down physiological dormancy, and involves the addition of moisture to the seeds so they absorb water, and they are then subjected to a period of moist chilling to after-ripen the embryo. Sowing in late summer and fall and allowing to overwinter under cool conditions is an effective way to stratify seeds; some seeds respond more favorably to periods of oscillating temperatures which are a part of the natural environment. Leaching or the soaking in water removes chemical inhibitors in some seeds that prevent germination. Rain and melting snow naturally accomplish this task. For seeds planted in gardens, running water is best – if soaked in a container, 12 to 24 hours of soaking is sufficient. Soaking longer, especially in stagnant water, can result in oxygen starvation and seed death. Seeds with hard seed coats can be soaked in hot water to break open the impermeable cell layers that prevent water intake. Other methods used to assist in the germination of seeds that have dormancy include prechilling, predrying, daily alternation of temperature, light exposure, potassium nitrate, the use of plant growth regulators, such as gibberellins, cytokinins, ethylene, thiourea, sodium hypochlorite, and others. Some seeds germinate best after a fire. For some seeds, fire cracks hard seed coats, while in others, chemical dormancy is broken in reaction to the presence of smoke. Liquid smoke is often used by gardeners to assist in the germination of these species. Sterile seeds Seeds may be sterile for few reasons: they may have been irradiated, unpollinated, cells lived past expectancy, or bred for the purpose. Evolution and origin of seeds The issue of the origin of seed plants remains unsolved. However, more and more data tends to place this origin in the middle Devonian. The description in 2004 of the proto-seed Runcaria heinzelinii in the Givetian of Belgium is an indication of that ancient origin of seed-plants. As with modern ferns, most land plants before this time reproduced by sending into the air spores that would land and become whole new plants. Taxonomists have described early "true" seeds from the upper Devonian, which probably became the theater of their true first evolutionary radiation. With this radiation came an evolution of seed size, shape, dispersal and eventually the radiation of gymnosperms and angiosperms and monocotyledons and dicotyledons. Seed plants progressively became one of the major elements of nearly all ecosystems. True to the seed Also called growing true, refers to plants whose seed will yield the same type of plant as the original plant. Open pollinated plants, which include heirlooms, will almost always grow true to seed if another variety does not cross-pollinate them. Seed microbiome Seeds harbor a diverse microbial community. Most of these microorganisms are transmitted from the seed to the developing seedlings. Economic importance Seed market In the United States farmers spent $22 billion on seeds in 2018, a 35 percent increase since 2010. DowDuPont and Monsanto account for 72 percent of corn and soybean seed sales in the U.S. with the average price of a bag of GMO corn seed is priced at $270. Seed production Seed production in natural plant populations varies widely from year to year in response to weather variables, insects and diseases, and internal cycles within the plants themselves. Over a 20-year period, for example, forests composed of loblolly pine and shortleaf pine produced from 0 to nearly 5.5 million sound pine seeds per hectare. Over this period, there were six bumper, five poor, and nine good seed crops, when evaluated for production of adequate seedlings for natural forest reproduction. Edible seeds Many seeds are edible and the majority of human calories comes from seeds, especially from cereals, legumes and nuts. Seeds also provide most cooking oils, many beverages and spices and some important food additives. In different seeds the seed embryo or the endosperm dominates and provides most of the nutrients. The storage proteins of the embryo and endosperm differ in their amino acid content and physical properties. For example, the gluten of wheat, important in providing the elastic property to bread dough is strictly an endosperm protein. Seeds are used to propagate many crops such as cereals, legumes, forest trees, turfgrasses, and pasture grasses. Particularly in developing countries, a major constraint faced is the inadequacy of the marketing channels to get the seed to poor farmers. Thus the use of farmer-retained seed remains quite common. Seeds are also eaten by animals (seed predation), and are also fed to livestock or provided as birdseed. Poison and food safety While some seeds are edible, others are harmful, poisonous or deadly. Plants and seeds often contain chemical compounds to discourage herbivores and seed predators. In some cases, these compounds simply taste bad (such as in mustard), but other compounds are toxic or break down into toxic compounds within the digestive system. Children, being smaller than adults, are more susceptible to poisoning by plants and seeds. A deadly poison, ricin, comes from seeds of the castor bean. Reported lethal doses are anywhere from two to eight seeds, though only a few deaths have been reported when castor beans have been ingested by animals. In addition, seeds containing amygdalin – apple, apricot, bitter almond, peach, plum, cherry, quince, and others – when consumed in sufficient amounts, may cause cyanide poisoning. Other seeds that contain poisons include annona, cotton, custard apple, datura, uncooked durian, golden chain, horse-chestnut, larkspur, locoweed, lychee, nectarine, rambutan, rosary pea, sour sop, sugar apple, wisteria, and yew. The seeds of the strychnine tree are also poisonous, containing the poison strychnine. The seeds of many legumes, including the common bean (Phaseolus vulgaris), contain proteins called lectins which can cause gastric distress if the beans are eaten without cooking. The common bean and many others, including the soybean, also contain trypsin inhibitors which interfere with the action of the digestive enzyme trypsin. Normal cooking processes degrade lectins and trypsin inhibitors to harmless forms. Other uses Cotton fiber grows attached to cotton plant seeds. Other seed fibers are from kapok and milkweed. Many important nonfood oils are extracted from seeds. Linseed oil is used in paints. Oil from jojoba and crambe are similar to whale oil. Seeds are the source of some medicines including castor oil, tea tree oil and the quack cancer drug Laetrile. Many seeds have been used as beads in necklaces and rosaries including Job's tears, Chinaberry, rosary pea, and castor bean. However, the latter three are also poisonous. Other seed uses include: Seeds once used as weights for balances. Seeds used as toys by children, such as for the game Conkers. Resin from Clusia rosea seeds used to caulk boats. Nematicide from milkweed seeds. Cottonseed meal used as animal feed and fertilizer. Seed records The oldest viable carbon-14-dated seed that has grown into a plant was a Judean date palm seed about 2,000 years old, recovered from excavations at Herod the Great's palace on Masada in Israel. It was germinated in 2005. (A reported regeneration of Silene stenophylla (narrow-leafed campion) from material preserved for 31,800 years in the Siberian permafrost was achieved using fruit tissue, not seed.) The largest seed is produced by the coco de mer, or "double coconut palm", Lodoicea maldivica. The entire fruit may weigh up to 23 kilograms (50 pounds) and usually contains a single seed. The smallest seeds are produced by epiphytic orchids. They are only 85 micrometers long, and weigh 0.81 micrograms. They have no endosperm and contain underdeveloped embryos. The earliest fossil seeds are around 365 million years old from the Late Devonian of West Virginia. The seeds are preserved immature ovules of the plant Elkinsia polymorpha. In religion The Book of Genesis in the Old Testament begins with an explanation of how all plant forms began: And God said, Let the earth bring forth grass, the herb yielding seed, and the fruit tree yielding fruit after his kind, whose seed is in itself, upon the earth: and it was so. And the earth brought forth grass, and herb yielding seed after its kind, and the tree yielding fruit, whose seed was in itself, after its kind: and God saw that it was good. And the evening and the morning were the third day. The Quran speaks of seed germination thus: It is Allah Who causeth the seed-grain and the date-stone to split and sprout. He causeth the living to issue from the dead, and He is the one to cause the dead to issue from the living. That is Allah: then how are ye deluded away from the truth?	Biology and health sciences	Biology	null
37710	https://en.wikipedia.org/wiki/Enthalpy%20of%20vaporization	Enthalpy of vaporization	In thermodynamics, the enthalpy of vaporization (symbol ), also known as the (latent) heat of vaporization or heat of evaporation, is the amount of energy (enthalpy) that must be added to a liquid substance to transform a quantity of that substance into a gas. The enthalpy of vaporization is a function of the pressure and temperature at which the transformation (vaporization or evaporation) takes place. The enthalpy of vaporization is often quoted for the normal boiling temperature of the substance. Although tabulated values are usually corrected to 298 K, that correction is often smaller than the uncertainty in the measured value. The heat of vaporization is temperature-dependent, though a constant heat of vaporization can be assumed for small temperature ranges and for reduced temperature . The heat of vaporization diminishes with increasing temperature and it vanishes completely at a certain point called the critical temperature (). Above the critical temperature, the liquid and vapor phases are indistinguishable, and the substance is called a supercritical fluid. Units Values are usually quoted in J/mol, or kJ/mol (molar enthalpy of vaporization), although kJ/kg, or J/g (specific heat of vaporization), and older units like kcal/mol, cal/g and Btu/lb are sometimes still used among others. Enthalpy of condensation The enthalpy of condensation (or heat of condensation) is by definition equal to the enthalpy of vaporization with the opposite sign: enthalpy changes of vaporization are always positive (heat is absorbed by the substance), whereas enthalpy changes of condensation are always negative (heat is released by the substance). Thermodynamic background The enthalpy of vaporization can be written as It is equal to the increased internal energy of the vapor phase compared with the liquid phase, plus the work done against ambient pressure. The increase in the internal energy can be viewed as the energy required to overcome the intermolecular interactions in the liquid (or solid, in the case of sublimation). Hence helium has a particularly low enthalpy of vaporization, 0.0845 kJ/mol, as the van der Waals forces between helium atoms are particularly weak. On the other hand, the molecules in liquid water are held together by relatively strong hydrogen bonds, and its enthalpy of vaporization, 40.65 kJ/mol, is more than five times the energy required to heat the same quantity of water from 0 °C to 100 °C (cp = 75.3 J/K·mol). Care must be taken, however, when using enthalpies of vaporization to measure the strength of intermolecular forces, as these forces may persist to an extent in the gas phase (as is the case with hydrogen fluoride), and so the calculated value of the bond strength will be too low. This is particularly true of metals, which often form covalently bonded molecules in the gas phase: in these cases, the enthalpy of atomization must be used to obtain a true value of the bond energy. An alternative description is to view the enthalpy of condensation as the heat which must be released to the surroundings to compensate for the drop in entropy when a gas condenses to a liquid. As the liquid and gas are in equilibrium at the boiling point (Tb), ΔvG = 0, which leads to: As neither entropy nor enthalpy vary greatly with temperature, it is normal to use the tabulated standard values without any correction for the difference in temperature from 298 K. A correction must be made if the pressure is different from 100 kPa, as the entropy of an ideal gas is proportional to the logarithm of its pressure. The entropies of liquids vary little with pressure, as the coefficient of thermal expansion of a liquid is small. These two definitions are equivalent: the boiling point is the temperature at which the increased entropy of the gas phase overcomes the intermolecular forces. As a given quantity of matter always has a higher entropy in the gas phase than in a condensed phase ( is always positive), and from , the Gibbs free energy change falls with increasing temperature: gases are favored at higher temperatures, as is observed in practice. Vaporization enthalpy of electrolyte solutions Estimation of the enthalpy of vaporization of electrolyte solutions can be simply carried out using equations based on the chemical thermodynamic models, such as Pitzer model or TCPC model. Selected values Elements The vaporization of metals is a key step in metal vapor synthesis, which exploits the increased reactivity of metal atoms or small particles relative to the bulk elements. Other common substances Enthalpies of vaporization of common substances, measured at their respective standard boiling points:	Physical sciences	Thermodynamics	Physics
37738	https://en.wikipedia.org/wiki/Soil	Soil	Soil, also commonly referred to as earth, is a mixture of organic matter, minerals, gases, liquids, and organisms that together support the life of plants and soil organisms. Some scientific definitions distinguish dirt from soil by restricting the former term specifically to displaced soil. Soil consists of a solid phase of minerals and organic matter (the soil matrix), as well as a porous phase that holds gases (the soil atmosphere) and water (the soil solution). Accordingly, soil is a three-state system of solids, liquids, and gases. Soil is a product of several factors: the influence of climate, relief (elevation, orientation, and slope of terrain), organisms, and the soil's parent materials (original minerals) interacting over time. It continually undergoes development by way of numerous physical, chemical and biological processes, which include weathering with associated erosion. Given its complexity and strong internal connectedness, soil ecologists regard soil as an ecosystem. Most soils have a dry bulk density (density of soil taking into account voids when dry) between 1.1 and 1.6 g/cm3, though the soil particle density is much higher, in the range of 2.6 to 2.7 g/cm3. Little of the soil of planet Earth is older than the Pleistocene and none is older than the Cenozoic, although fossilized soils are preserved from as far back as the Archean. Collectively the Earth's body of soil is called the pedosphere. The pedosphere interfaces with the lithosphere, the hydrosphere, the atmosphere, and the biosphere. Soil has four important functions: as a medium for plant growth as a means of water storage, supply, and purification as a modifier of Earth's atmosphere as a habitat for organisms All of these functions, in their turn, modify the soil and its properties. Soil science has two basic branches of study: edaphology and pedology. Edaphology studies the influence of soils on living things. Pedology focuses on the formation, description (morphology), and classification of soils in their natural environment. In engineering terms, soil is included in the broader concept of regolith, which also includes other loose material that lies above the bedrock, as can be found on the Moon and other celestial objects. Processes Soil is a major component of the Earth's ecosystem. The world's ecosystems are impacted in far-reaching ways by the processes carried out in the soil, with effects ranging from ozone depletion and global warming to rainforest destruction and water pollution. With respect to Earth's carbon cycle, soil acts as an important carbon reservoir, and it is potentially one of the most reactive to human disturbance and climate change. As the planet warms, it has been predicted that soils will add carbon dioxide to the atmosphere due to increased biological activity at higher temperatures, a positive feedback (amplification). This prediction has, however, been questioned on consideration of more recent knowledge on soil carbon turnover. Soil acts as an engineering medium, a habitat for soil organisms, a recycling system for nutrients and organic wastes, a regulator of water quality, a modifier of atmospheric composition, and a medium for plant growth, making it a critically important provider of ecosystem services. Since soil has a tremendous range of available niches and habitats, it contains a prominent part of the Earth's genetic diversity. A gram of soil can contain billions of organisms, belonging to thousands of species, mostly microbial and largely still unexplored. Soil has a mean prokaryotic density of roughly 108 organisms per gram, whereas the ocean has no more than 107 prokaryotic organisms per milliliter (gram) of seawater. Organic carbon held in soil is eventually returned to the atmosphere through the process of respiration carried out by heterotrophic organisms, but a substantial part is retained in the soil in the form of soil organic matter; tillage usually increases the rate of soil respiration, leading to the depletion of soil organic matter. Since plant roots need oxygen, aeration is an important characteristic of soil. This ventilation can be accomplished via networks of interconnected soil pores, which also absorb and hold rainwater making it readily available for uptake by plants. Since plants require a nearly continuous supply of water, but most regions receive sporadic rainfall, the water-holding capacity of soils is vital for plant survival. Soils can effectively remove impurities, kill disease agents, and degrade contaminants, this latter property being called natural attenuation. Typically, soils maintain a net absorption of oxygen and methane and undergo a net release of carbon dioxide and nitrous oxide. Soils offer plants physical support, air, water, temperature moderation, nutrients, and protection from toxins. Soils provide readily available nutrients to plants and animals by converting dead organic matter into various nutrient forms. Composition A typical soil is about 50% solids (45% mineral and 5% organic matter), and 50% voids (or pores) of which half is occupied by water and half by gas. The percent soil mineral and organic content can be treated as a constant (in the short term), while the percent soil water and gas content is considered highly variable whereby a rise in one is simultaneously balanced by a reduction in the other. The pore space allows for the infiltration and movement of air and water, both of which are critical for life existing in soil. Compaction, a common problem with soils, reduces this space, preventing air and water from reaching plant roots and soil organisms. Given sufficient time, an undifferentiated soil will evolve a soil profile that consists of two or more layers, referred to as soil horizons. These differ in one or more properties such as in their texture, structure, density, porosity, consistency, temperature, color, and reactivity. The horizons differ greatly in thickness and generally lack sharp boundaries; their development is dependent on the type of parent material, the processes that modify those parent materials, and the soil-forming factors that influence those processes. The biological influences on soil properties are strongest near the surface, while the geochemical influences on soil properties increase with depth. Mature soil profiles typically include three basic master horizons: A, B, and C. The solum normally includes the A and B horizons. The living component of the soil is largely confined to the solum, and is generally more prominent in the A horizon. It has been suggested that the pedon, a column of soil extending vertically from the surface to the underlying parent material and large enough to show the characteristics of all its horizons, could be subdivided in the humipedon (the living part, where most soil organisms are dwelling, corresponding to the humus form), the copedon (in intermediary position, where most weathering of minerals takes place) and the lithopedon (in contact with the subsoil). The soil texture is determined by the relative proportions of the individual particles of sand, silt, and clay that make up the soil. The interaction of the individual mineral particles with organic matter, water, gases via biotic and abiotic processes causes those particles to flocculate (stick together) to form aggregates or peds. Where these aggregates can be identified, a soil can be said to be developed, and can be described further in terms of color, porosity, consistency, reaction (acidity), etc. Water is a critical agent in soil development due to its involvement in the dissolution, precipitation, erosion, transport, and deposition of the materials of which a soil is composed. The mixture of water and dissolved or suspended materials that occupy the soil pore space is called the soil solution. Since soil water is never pure water, but contains hundreds of dissolved organic and mineral substances, it may be more accurately called the soil solution. Water is central to the dissolution, precipitation and leaching of minerals from the soil profile. Finally, water affects the type of vegetation that grows in a soil, which in turn affects the development of the soil, a complex feedback which is exemplified in the dynamics of banded vegetation patterns in semi-arid regions. Soils supply plants with nutrients, most of which are held in place by particles of clay and organic matter (colloids) The nutrients may be adsorbed on clay mineral surfaces, bound within clay minerals (absorbed), or bound within organic compounds as part of the living organisms or dead soil organic matter. These bound nutrients interact with soil water to buffer the soil solution composition (attenuate changes in the soil solution) as soils wet up or dry out, as plants take up nutrients, as salts are leached, or as acids or alkalis are added. Plant nutrient availability is affected by soil pH, which is a measure of the hydrogen ion activity in the soil solution. Soil pH is a function of many soil forming factors, and is generally lower (more acidic) where weathering is more advanced. Most plant nutrients, with the exception of nitrogen, originate from the minerals that make up the soil parent material. Some nitrogen originates from rain as dilute nitric acid and ammonia, but most of the nitrogen is available in soils as a result of nitrogen fixation by bacteria. Once in the soil-plant system, most nutrients are recycled through living organisms, plant and microbial residues (soil organic matter), mineral-bound forms, and the soil solution. Both living soil organisms (microbes, animals and plant roots) and soil organic matter are of critical importance to this recycling, and thereby to soil formation and soil fertility. Microbial soil enzymes may release nutrients from minerals or organic matter for use by plants and other microorganisms, sequester (incorporate) them into living cells, or cause their loss from the soil by volatilisation (loss to the atmosphere as gases) or leaching. Formation Soil is said to be formed when organic matter has accumulated and colloids are washed downward, leaving deposits of clay, humus, iron oxide, carbonate, and gypsum, producing a distinct layer called the B horizon. This is a somewhat arbitrary definition as mixtures of sand, silt, clay and humus will support biological and agricultural activity before that time. These constituents are moved from one level to another by water and animal activity. As a result, layers (horizons) form in the soil profile. The alteration and movement of materials within a soil causes the formation of distinctive soil horizons. However, more recent definitions of soil embrace soils without any organic matter, such as those regoliths that formed on Mars and analogous conditions in planet Earth deserts. An example of the development of a soil would begin with the weathering of lava flow bedrock, which would produce the purely mineral-based parent material from which the soil texture forms. Soil development would proceed most rapidly from bare rock of recent flows in a warm climate, under heavy and frequent rainfall. Under such conditions, plants (in a first stage nitrogen-fixing lichens and cyanobacteria then epilithic higher plants) become established very quickly on basaltic lava, even though there is very little organic material. Basaltic minerals commonly weather relatively quickly, according to the Goldich dissolution series. The plants are supported by the porous rock as it is filled with nutrient-bearing water that carries minerals dissolved from the rocks. Crevasses and pockets, local topography of the rocks, would hold fine materials and harbour plant roots. The developing plant roots are associated with mineral-weathering mycorrhizal fungi that assist in breaking up the porous lava, and by these means organic matter and a finer mineral soil accumulate with time. Such initial stages of soil development have been described on volcanoes, inselbergs, and glacial moraines. How soil formation proceeds is influenced by at least five classic factors that are intertwined in the evolution of a soil: parent material, climate, topography (relief), organisms, and time. When reordered to climate, relief, organisms, parent material, and time, they form the acronym CROPT. Physical properties The physical properties of soils, in order of decreasing importance for ecosystem services such as crop production, are texture, structure, bulk density, porosity, consistency, temperature, colour and resistivity. Soil texture is determined by the relative proportion of the three kinds of soil mineral particles, called soil separates: sand, silt, and clay. At the next larger scale, soil structures called peds or more commonly soil aggregates are created from the soil separates when iron oxides, carbonates, clay, silica and humus, coat particles and cause them to adhere into larger, relatively stable secondary structures. Soil bulk density, when determined at standardized moisture conditions, is an estimate of soil compaction. Soil porosity consists of the void part of the soil volume and is occupied by gases or water. Soil consistency is the ability of soil materials to stick together. Soil temperature and colour are self-defining. Resistivity refers to the resistance to conduction of electric currents and affects the rate of corrosion of metal and concrete structures which are buried in soil. These properties vary through the depth of a soil profile, i.e. through soil horizons. Most of these properties determine the aeration of the soil and the ability of water to infiltrate and to be held within the soil. Soil moisture Soil water content can be measured as volume or weight. Soil moisture levels, in order of decreasing water content, are saturation, field capacity, wilting point, air dry, and oven dry. Field capacity describes a drained wet soil at the point water content reaches equilibrium with gravity. Irrigating soil above field capacity risks percolation losses. Wilting point describes the dry limit for growing plants. During growing season, soil moisture is unaffected by functional groups or specie richness. Available water capacity is the amount of water held in a soil profile available to plants. As water content drops, plants have to work against increasing forces of adhesion and sorptivity to withdraw water. Irrigation scheduling avoids moisture stress by replenishing depleted water before stress is induced. Capillary action is responsible for moving groundwater from wet regions of the soil to dry areas. Subirrigation designs (e.g., wicking beds, sub-irrigated planters) rely on capillarity to supply water to plant roots. Capillary action can result in an evaporative concentration of salts, causing land degradation through salination. Soil moisture measurement—measuring the water content of the soil, as can be expressed in terms of volume or weight—can be based on in situ probes (e.g., capacitance probes, neutron probes), or remote sensing methods. Soil moisture measurement is an important factor in determining changes in soil activity. Soil gas The atmosphere of soil, or soil gas, is very different from the atmosphere above. The consumption of oxygen by microbes and plant roots, and their release of carbon dioxide, decreases oxygen and increases carbon dioxide concentration. Atmospheric CO2 concentration is 0.04%, but in the soil pore space it may range from 10 to 100 times that level, thus potentially contributing to the inhibition of root respiration. Calcareous soils regulate CO2 concentration by carbonate buffering, contrary to acid soils in which all CO2 respired accumulates in the soil pore system. At extreme levels, CO2 is toxic. This suggests a possible negative feedback control of soil CO2 concentration through its inhibitory effects on root and microbial respiration (also called soil respiration). In addition, the soil voids are saturated with water vapour, at least until the point of maximal hygroscopicity, beyond which a vapour-pressure deficit occurs in the soil pore space. Adequate porosity is necessary, not just to allow the penetration of water, but also to allow gases to diffuse in and out. Movement of gases is by diffusion from high concentrations to lower, the diffusion coefficient decreasing with soil compaction. Oxygen from above atmosphere diffuses in the soil where it is consumed and levels of carbon dioxide in excess of above atmosphere diffuse out with other gases (including greenhouse gases) as well as water. Soil texture and structure strongly affect soil porosity and gas diffusion. It is the total pore space (porosity) of soil, not the pore size, and the degree of pore interconnection (or conversely pore sealing), together with water content, air turbulence and temperature, that determine the rate of diffusion of gases into and out of soil. Platy soil structure and soil compaction (low porosity) impede gas flow, and a deficiency of oxygen may encourage anaerobic bacteria to reduce (strip oxygen) from nitrate NO3 to the gases N2, N2O, and NO, which are then lost to the atmosphere, thereby depleting the soil of nitrogen, a detrimental process called denitrification. Aerated soil is also a net sink of methane (CH4) but a net producer of methane (a strong heat-absorbing greenhouse gas) when soils are depleted of oxygen and subject to elevated temperatures. Soil atmosphere is also the seat of emissions of volatiles other than carbon and nitrogen oxides from various soil organisms, e.g. roots, bacteria, fungi, animals. These volatiles are used as chemical cues, making soil atmosphere the seat of interaction networks playing a decisive role in the stability, dynamics and evolution of soil ecosystems. Biogenic soil volatile organic compounds are exchanged with the aboveground atmosphere, in which they are just 1–2 orders of magnitude lower than those from aboveground vegetation. Humans can get some idea of the soil atmosphere through the well-known 'after-the-rain' scent, when infiltering rainwater flushes out the whole soil atmosphere after a drought period, or when soil is excavated, a bulk property attributed in a reductionist manner to particular biochemical compounds such as petrichor or geosmin. Solid phase (soil matrix) Soil particles can be classified by their chemical composition (mineralogy) as well as their size. The particle size distribution of a soil, its texture, determines many of the properties of that soil, in particular hydraulic conductivity and water potential, but the mineralogy of those particles can strongly modify those properties. The mineralogy of the finest soil particles, clay, is especially important. Soil biodiversity Large numbers of microbes, animals, plants and fungi are living in soil. However, biodiversity in soil is much harder to study as most of this life is invisible, hence estimates about soil biodiversity have been unsatisfactory. A recent study suggested that soil is likely home to 59 ± 15% of the species on Earth. Enchytraeidae (worms) have the greatest percentage of species in soil (98.6%), followed by fungi (90%), plants (85.5%), and termites (Isoptera) (84.2%). Many other groups of animals have substantial fractions of species living in soil, e.g. about 30% of insects, and close to 50% of arachnids. While most vertebrates live above ground (ignoring aquatic species), many species are fossorial, that is, they live in soil, such as most blind snakes. Chemistry The chemistry of a soil determines its ability to supply available plant nutrients and affects its physical properties and the health of its living population. In addition, a soil's chemistry also determines its corrosivity, stability, and ability to absorb pollutants and to filter water. It is the surface chemistry of mineral and organic colloids that determines soil's chemical properties. A colloid is a small, insoluble particle ranging in size from 1 nanometer to 1 micrometer, thus small enough to remain suspended by Brownian motion in a fluid medium without settling. Most soils contain organic colloidal particles called humus as well as the inorganic colloidal particles of clays. The very high specific surface area of colloids and their net electrical charges give soil its ability to hold and release ions. Negatively charged sites on colloids attract and release cations in what is referred to as cation exchange. Cation-exchange capacity is the amount of exchangeable cations per unit weight of dry soil and is expressed in terms of milliequivalents of positively charged ions per 100 grams of soil (or centimoles of positive charge per kilogram of soil; cmolc/kg). Similarly, positively charged sites on colloids can attract and release anions in the soil, giving the soil anion exchange capacity. Cation and anion exchange The cation exchange, that takes place between colloids and soil water, buffers (moderates) soil pH, alters soil structure, and purifies percolating water by adsorbing cations of all types, both useful and harmful. The negative or positive charges on colloid particles make them able to hold cations or anions, respectively, to their surfaces. The charges result from four sources. Isomorphous substitution occurs in clay during its formation, when lower-valence cations substitute for higher-valence cations in the crystal structure. Substitutions in the outermost layers are more effective than for the innermost layers, as the electric charge strength drops off as the square of the distance. The net result is oxygen atoms with net negative charge and the ability to attract cations. Edge-of-clay oxygen atoms are not in balance ionically as the tetrahedral and octahedral structures are incomplete. Hydroxyls may substitute for oxygens of the silica layers, a process called hydroxylation. When the hydrogens of the clay hydroxyls are ionised into solution, they leave the oxygen with a negative charge (anionic clays). Hydrogens of humus hydroxyl groups may also be ionised into solution, leaving, similarly to clay, an oxygen with a negative charge. Cations held to the negatively charged colloids resist being washed downward by water and are out of reach of plant roots, thereby preserving the soil fertility in areas of moderate rainfall and low temperatures. There is a hierarchy in the process of cation exchange on colloids, as cations differ in the strength of adsorption by the colloid and hence their ability to replace one another (ion exchange). If present in equal amounts in the soil water solution: Al3+ replaces H+ replaces Ca2+ replaces Mg2+ replaces K+ same as replaces Na+ If one cation is added in large amounts, it may replace the others by the sheer force of its numbers. This is called law of mass action. This is largely what occurs with the addition of cationic fertilisers (potash, lime). As the soil solution becomes more acidic (low pH, meaning an abundance of H+), the other cations more weakly bound to colloids are pushed into solution as hydrogen ions occupy exchange sites (protonation). A low pH may cause the hydrogen of hydroxyl groups to be pulled into solution, leaving charged sites on the colloid available to be occupied by other cations. This ionisation of hydroxy groups on the surface of soil colloids creates what is described as pH-dependent surface charges. Unlike permanent charges developed by isomorphous substitution, pH-dependent charges are variable and increase with increasing pH. Freed cations can be made available to plants but are also prone to be leached from the soil, possibly making the soil less fertile. Plants are able to excrete H+ into the soil through the synthesis of organic acids and by that means, change the pH of the soil near the root and push cations off the colloids, thus making those available to the plant. Cation exchange capacity (CEC) Cation exchange capacity is the soil's ability to remove cations from the soil water solution and sequester those to be exchanged later as the plant roots release hydrogen ions to the solution. CEC is the amount of exchangeable hydrogen cation (H+) that will combine with 100 grams dry weight of soil and whose measure is one milliequivalents per 100 grams of soil (1 meq/100 g). Hydrogen ions have a single charge and one-thousandth of a gram of hydrogen ions per 100 grams dry soil gives a measure of one milliequivalent of hydrogen ion. Calcium, with an atomic weight 40 times that of hydrogen and with a valence of two, converts to = 20 milliequivalents of hydrogen ion per 100 grams of dry soil or 20 meq/100 g. The modern measure of CEC is expressed as centimoles of positive charge per kilogram (cmol/kg) of oven-dry soil. Most of the soil's CEC occurs on clay and humus colloids, and the lack of those in hot, humid, wet climates (such as tropical rainforests), due to leaching and decomposition, respectively, explains the apparent sterility of tropical soils. Live plant roots also have some CEC, linked to their specific surface area. Anion exchange capacity (AEC) Anion exchange capacity is the soil's ability to remove anions (such as nitrate, phosphate) from the soil water solution and sequester those for later exchange as the plant roots release carbonate anions to the soil water solution. Those colloids which have low CEC tend to have some AEC. Amorphous and sesquioxide clays have the highest AEC, followed by the iron oxides. Levels of AEC are much lower than for CEC, because of the generally higher rate of positively (versus negatively) charged surfaces on soil colloids, to the exception of variable-charge soils. Phosphates tend to be held at anion exchange sites. Iron and aluminum hydroxide clays are able to exchange their hydroxide anions (OH−) for other anions. The order reflecting the strength of anion adhesion is as follows: replaces replaces replaces Cl− The amount of exchangeable anions is of a magnitude of tenths to a few milliequivalents per 100 g dry soil. As pH rises, there are relatively more hydroxyls, which will displace anions from the colloids and force them into solution and out of storage; hence AEC decreases with increasing pH (alkalinity). Reactivity (pH) Soil reactivity is expressed in terms of pH and is a measure of the acidity or alkalinity of the soil. More precisely, it is a measure of hydronium concentration in an aqueous solution and ranges in values from 0 to 14 (acidic to basic) but practically speaking for soils, pH ranges from 3.5 to 9.5, as pH values beyond those extremes are toxic to life forms. At 25 °C an aqueous solution that has a pH of 3.5 has 10−3.5 moles H3O+ (hydronium ions) per litre of solution (and also 10−10.5 moles per litre OH−). A pH of 7, defined as neutral, has 10−7 moles of hydronium ions per litre of solution and also 10−7 moles of OH− per litre; since the two concentrations are equal, they are said to neutralise each other. A pH of 9.5 has 10−9.5 moles hydronium ions per litre of solution (and also 10−2.5 moles per litre OH−). A pH of 3.5 has one million times more hydronium ions per litre than a solution with pH of 9.5 ( or 106) and is more acidic. The effect of pH on a soil is to remove from the soil or to make available certain ions. Soils with high acidity tend to have toxic amounts of aluminium and manganese. As a result of a trade-off between toxicity and requirement most nutrients are better available to plants at moderate pH, although most minerals are more soluble in acid soils. Soil organisms are hindered by high acidity, and most agricultural crops do best with mineral soils of pH 6.5 and organic soils of pH 5.5. Given that at low pH toxic metals (e.g. cadmium, zinc, lead) are positively charged as cations and organic pollutants are in non-ionic form, thus both made more available to organisms, it has been suggested that plants, animals and microbes commonly living in acid soils are pre-adapted to every kind of pollution, whether of natural or human origin. In high rainfall areas, soils tend to acidify as the basic cations are forced off the soil colloids by the mass action of hydronium ions from usual or unusual rain acidity against those attached to the colloids. High rainfall rates can then wash the nutrients out, leaving the soil inhabited only by those organisms which are particularly efficient to uptake nutrients in very acid conditions, like in tropical rainforests. Once the colloids are saturated with H3O+, the addition of any more hydronium ions or aluminum hydroxyl cations drives the pH even lower (more acidic) as the soil has been left with no buffering capacity. In areas of extreme rainfall and high temperatures, the clay and humus may be washed out, further reducing the buffering capacity of the soil. In low rainfall areas, unleached calcium pushes pH to 8.5 and with the addition of exchangeable sodium, soils may reach pH 10. Beyond a pH of 9, plant growth is reduced. High pH results in low micro-nutrient mobility, but water-soluble chelates of those nutrients can correct the deficit. Sodium can be reduced by the addition of gypsum (calcium sulphate) as calcium adheres to clay more tightly than does sodium causing sodium to be pushed into the soil water solution where it can be washed out by an abundance of water. Base saturation percentage There are acid-forming cations (e.g. hydronium, aluminium, iron) and there are base-forming cations (e.g. calcium, magnesium, sodium). The fraction of the negatively-charged soil colloid exchange sites (CEC) that are occupied by base-forming cations is called base saturation. If a soil has a CEC of 20 meq and 5 meq are aluminium and hydronium cations (acid-forming), the remainder of positions on the colloids () are assumed occupied by base-forming cations, so that the base saturation is (the compliment 25% is assumed acid-forming cations). Base saturation is almost in direct proportion to pH (it increases with increasing pH). It is of use in calculating the amount of lime needed to neutralise an acid soil (lime requirement). The amount of lime needed to neutralize a soil must take account of the amount of acid forming ions on the colloids (exchangeable acidity), not just those in the soil water solution (free acidity). The addition of enough lime to neutralize the soil water solution will be insufficient to change the pH, as the acid forming cations stored on the soil colloids will tend to restore the original pH condition as they are pushed off those colloids by the calcium of the added lime. Buffering The resistance of soil to change in pH, as a result of the addition of acid or basic material, is a measure of the buffering capacity of a soil and (for a particular soil type) increases as the CEC increases. Hence, pure sand has almost no buffering ability, though soils high in colloids (whether mineral or organic) have high buffering capacity. Buffering occurs by cation exchange and neutralisation. However, colloids are not the only regulators of soil pH. The role of carbonates should be underlined, too. More generally, according to pH levels, several buffer systems take precedence over each other, from calcium carbonate buffer range to iron buffer range. The addition of a small amount of highly basic aqueous ammonia to a soil will cause the ammonium to displace hydronium ions from the colloids, and the end product is water and colloidally fixed ammonium, but little permanent change overall in soil pH. The addition of a small amount of lime, Ca(OH)2, will displace hydronium ions from the soil colloids, causing the fixation of calcium to colloids and the evolution of CO2 and water, with little permanent change in soil pH. The above are examples of the buffering of soil pH. The general principal is that an increase in a particular cation in the soil water solution will cause that cation to be fixed to colloids (buffered) and a decrease in solution of that cation will cause it to be withdrawn from the colloid and moved into solution (buffered). The degree of buffering is often related to the CEC of the soil; the greater the CEC, the greater the buffering capacity of the soil. Redox Soil chemical reactions involve some combination of proton and electron transfer. Oxidation occurs if there is a loss of electrons in the transfer process while reduction occurs if there is a gain of electrons. Reduction potential is measured in volts or millivolts. Soil microbial communities develop along electron transport chains, forming electrically conductive biofilms, and developing networks of bacterial nanowires. Redox factors in soil development, where formation of redoximorphic color features provides critical information for soil interpretation. Understanding the redox gradient is important to managing carbon sequestration, bioremediation, wetland delineation, and soil-based microbial fuel cells. Nutrients Seventeen elements or nutrients are essential for plant growth and reproduction. They are carbon (C), hydrogen (H), oxygen (O), nitrogen (N), phosphorus (P), potassium (K), sulfur (S), calcium (Ca), magnesium (Mg), iron (Fe), boron (B), manganese (Mn), copper (Cu), zinc (Zn), molybdenum (Mo), nickel (Ni) and chlorine (Cl). Nutrients required for plants to complete their life cycle are considered essential nutrients. Nutrients that enhance the growth of plants but are not necessary to complete the plant's life cycle are considered non-essential. With the exception of carbon, hydrogen and oxygen, which are supplied by carbon dioxide and water, and nitrogen, provided through nitrogen fixation, the nutrients derive originally from the mineral component of the soil. The Law of the Minimum expresses that when the available form of a nutrient is not in enough proportion in the soil solution, then other nutrients cannot be taken up at an optimum rate by a plant. A particular nutrient ratio of the soil solution is thus mandatory for optimizing plant growth, a value which might differ from nutrient ratios calculated from plant composition. Plant uptake of nutrients can only proceed when they are present in a plant-available form. In most situations, nutrients are absorbed in an ionic form from (or together with) soil water. Although minerals are the origin of most nutrients, and the bulk of most nutrient elements in the soil is held in crystalline form within primary and secondary minerals, they weather too slowly to support rapid plant growth. For example, the application of finely ground minerals, feldspar and apatite, to soil seldom provides the necessary amounts of potassium and phosphorus at a rate sufficient for good plant growth, as most of the nutrients remain bound in the crystals of those minerals. The nutrients adsorbed onto the surfaces of clay colloids and soil organic matter provide a more accessible reservoir of many plant nutrients (e.g. K, Ca, Mg, P, Zn). As plants absorb the nutrients from the soil water, the soluble pool is replenished from the surface-bound pool. The decomposition of soil organic matter by microorganisms is another mechanism whereby the soluble pool of nutrients is replenished – this is important for the supply of plant-available N, S, P, and B from soil. Gram for gram, the capacity of humus to hold nutrients and water is far greater than that of clay minerals, most of the soil cation exchange capacity arising from charged carboxylic groups on organic matter. However, despite the great capacity of humus to retain water once water-soaked, its high hydrophobicity decreases its wettability once dry. All in all, small amounts of humus may remarkably increase the soil's capacity to promote plant growth. Soil organic matter The organic material in soil is made up of organic compounds and includes plant, animal and microbial material, both living and dead. A typical soil has a biomass composition of 70% microorganisms, 22% macrofauna, and 8% roots. The living component of an acre of soil may include 900 lb of earthworms, 2400 lb of fungi, 1500 lb of bacteria, 133 lb of protozoa and 890 lb of arthropods and algae. A few percent of the soil organic matter, with small residence time, consists of the microbial biomass and metabolites of bacteria, molds, and actinomycetes that work to break down the dead organic matter. Were it not for the action of these micro-organisms, the entire carbon dioxide part of the atmosphere would be sequestered as organic matter in the soil. However, in the same time soil microbes contribute to carbon sequestration in the topsoil through the formation of stable humus. In the aim to sequester more carbon in the soil for alleviating the greenhouse effect it would be more efficient in the long-term to stimulate humification than to decrease litter decomposition. The main part of soil organic matter is a complex assemblage of small organic molecules, collectively called humus or humic substances. The use of these terms, which do not rely on a clear chemical classification, has been considered as obsolete. Other studies showed that the classical notion of molecule is not convenient for humus, which escaped most attempts done over two centuries to resolve it in unit components, but still is chemically distinct from polysaccharides, lignins and proteins. Most living things in soils, including plants, animals, bacteria, and fungi, are dependent on organic matter for nutrients and/or energy. Soils have organic compounds in varying degrees of decomposition, the rate of which is dependent on the temperature, soil moisture, and aeration. Bacteria and fungi feed on the raw organic matter, which are fed upon by protozoa, which in turn are fed upon by nematodes, annelids and arthropods, themselves able to consume and transform raw or humified organic matter. This has been called the soil food web, through which all organic matter is processed as in a digestive system. Organic matter holds soils open, allowing the infiltration of air and water, and may hold as much as twice its weight in water. Many soils, including desert and rocky-gravel soils, have little or no organic matter. Soils that are all organic matter, such as peat (histosols), are infertile. In its earliest stage of decomposition, the original organic material is often called raw organic matter. The final stage of decomposition is called humus. In grassland, much of the organic matter added to the soil is from the deep, fibrous, grass root systems. By contrast, tree leaves falling on the forest floor are the principal source of soil organic matter in the forest. Another difference is the frequent occurrence in the grasslands of fires that destroy large amounts of aboveground material but stimulate even greater contributions from roots. Also, the much greater acidity under any forests inhibits the action of certain soil organisms that otherwise would mix much of the surface litter into the mineral soil. As a result, the soils under grasslands generally develop a thicker A horizon with a deeper distribution of organic matter than in comparable soils under forests, which characteristically store most of their organic matter in the forest floor (O horizon) and thin A horizon. Humus Humus refers to organic matter that has been decomposed by soil microflora and fauna to the point where it is resistant to further breakdown. Humus usually constitutes only five percent of the soil or less by volume, but it is an essential source of nutrients and adds important textural qualities crucial to soil health and plant growth. Humus also feeds arthropods, termites and earthworms which further improve the soil. The end product, humus, is suspended in colloidal form in the soil solution and forms a weak acid that can attack silicate minerals by chelating their iron and aluminum atoms. Humus has a high cation and anion exchange capacity that on a dry weight basis is many times greater than that of clay colloids. It also acts as a buffer, like clay, against changes in pH and soil moisture. Humic acids and fulvic acids, which begin as raw organic matter, are important constituents of humus. After the death of plants, animals, and microbes, microbes begin to feed on the residues through their production of extra-cellular soil enzymes, resulting finally in the formation of humus. As the residues break down, only molecules made of aliphatic and aromatic hydrocarbons, assembled and stabilized by oxygen and hydrogen bonds, remain in the form of complex molecular assemblages collectively called humus. Humus is never pure in the soil, because it reacts with metals and clays to form complexes which further contribute to its stability and to soil structure. Although the structure of humus has in itself few nutrients (with the exception of constitutive metals such as calcium, iron and aluminum) it is able to attract and link, by weak bonds, cation and anion nutrients that can further be released into the soil solution in response to selective root uptake and changes in soil pH, a process of paramount importance for the maintenance of fertility in tropical soils. Lignin is resistant to breakdown and accumulates within the soil. It also reacts with proteins, which further increases its resistance to decomposition, including enzymatic decomposition by microbes. Fats and waxes from plant matter have still more resistance to decomposition and persist in soils for thousand years, hence their use as tracers of past vegetation in buried soil layers. Clay soils often have higher organic contents that persist longer than soils without clay as the organic molecules adhere to and are stabilised by the clay. Proteins normally decompose readily, to the exception of scleroproteins, but when bound to clay particles they become more resistant to decomposition. As for other proteins clay particles absorb the enzymes exuded by microbes, decreasing enzyme activity while protecting extracellular enzymes from degradation. The addition of organic matter to clay soils can render that organic matter and any added nutrients inaccessible to plants and microbes for many years. A study showed increased soil fertility following the addition of mature compost to a clay soil. High soil tannin content can cause nitrogen to be sequestered as resistant tannin-protein complexes. Humus formation is a process dependent on the amount of plant material added each year and the type of base soil. Both are affected by climate and the type of organisms present. Soils with humus can vary in nitrogen content but typically have 3 to 6 percent nitrogen. Raw organic matter, as a reserve of nitrogen and phosphorus, is a vital component affecting soil fertility. Humus also absorbs water, and expands and shrinks between dry and wet states to a higher extent than clay, increasing soil porosity. Humus is less stable than the soil's mineral constituents, as it is reduced by microbial decomposition, and over time its concentration diminishes without the addition of new organic matter. However, humus in its most stable forms may persist over centuries if not millennia. Charcoal is a source of highly stable humus, called black carbon, which had been used traditionally to improve the fertility of nutrient-poor tropical soils. This very ancient practice, as ascertained in the genesis of Amazonian dark earths, has been renewed and became popular under the name of biochar. It has been suggested that biochar could be used to sequester more carbon in the fight against the greenhouse effect. Climatological influence The production, accumulation and degradation of organic matter are greatly dependent on climate. For example, when a thawing event occurs, the flux of soil gases with atmospheric gases is significantly influenced. Temperature, soil moisture and topography are the major factors affecting the accumulation of organic matter in soils. Organic matter tends to accumulate under wet or cold conditions where decomposer activity is impeded by low temperature or excess moisture which results in anaerobic conditions. Conversely, excessive rain and high temperatures of tropical climates enables rapid decomposition of organic matter and leaching of plant nutrients. Forest ecosystems on these soils rely on efficient recycling of nutrients and plant matter by the living plant and microbial biomass to maintain their productivity, a process which is disturbed by human activities. Excessive slope, in particular in the presence of cultivation for the sake of agriculture, may encourage the erosion of the top layer of soil which holds most of the raw organic material that would otherwise eventually become humus. Plant residue Cellulose and hemicellulose undergo fast decomposition by fungi and bacteria, with a half-life of 12–18 days in a temperate climate. Brown rot fungi can decompose the cellulose and hemicellulose, leaving the lignin and phenolic compounds behind. Starch, which is an energy storage system for plants, undergoes fast decomposition by bacteria and fungi. Lignin consists of polymers composed of 500 to 600 units with a highly branched, amorphous structure, linked to cellulose, hemicellulose and pectin in plant cell walls. Lignin undergoes very slow decomposition, mainly by white rot fungi and actinomycetes; its half-life under temperate conditions is about six months. Horizons A horizontal layer of the soil, whose physical features, composition and age are distinct from those above and beneath, is referred to as a soil horizon. The naming of a horizon is based on the type of material of which it is composed. Those materials reflect the duration of specific processes of soil formation. They are labelled using a shorthand notation of letters and numbers which describe the horizon in terms of its colour, size, texture, structure, consistency, root quantity, pH, voids, boundary characteristics and presence of nodules or concretions. No soil profile has all the major horizons. Some, called entisols, may have only one horizon or are currently considered as having no horizon, in particular incipient soils from unreclaimed mining waste deposits, moraines, volcanic cones sand dunes or alluvial terraces. Upper soil horizons may be lacking in truncated soils following wind or water ablation, with concomitant downslope burying of soil horizons, a natural process aggravated by agricultural practices such as tillage. The growth of trees is another source of disturbance, creating a micro-scale heterogeneity which is still visible in soil horizons once trees have died. By passing from a horizon to another, from the top to the bottom of the soil profile, one goes back in time, with past events registered in soil horizons like in sediment layers. Sampling pollen, testate amoebae and plant remains in soil horizons may help to reveal environmental changes (e.g. climate change, land use change) which occurred in the course of soil formation. Soil horizons can be dated by several methods such as radiocarbon, using pieces of charcoal provided they are of enough size to escape pedoturbation by earthworm activity and other mechanical disturbances. Fossil soil horizons from paleosols can be found within sedimentary rock sequences, allowing the study of past environments. The exposure of parent material to favourable conditions produces mineral soils that are marginally suitable for plant growth, as is the case in eroded soils. The growth of vegetation results in the production of organic residues which fall on the ground as litter for plant aerial parts (leaf litter) or are directly produced belowground for subterranean plant organs (root litter), and then release dissolved organic matter. The remaining surficial organic layer, called the O horizon, produces a more active soil due to the effect of the organisms that live within it. Organisms colonise and break down organic materials, making available nutrients upon which other plants and animals can live. After sufficient time, humus moves downward and is deposited in a distinctive organic-mineral surface layer called the A horizon, in which organic matter is mixed with mineral matter through the activity of burrowing animals, a process called pedoturbation. This natural process does not go to completion in the presence of conditions detrimental to soil life such as strong acidity, cold climate or pollution, stemming in the accumulation of undecomposed organic matter within a single organic horizon overlying the mineral soil and in the juxtaposition of humified organic matter and mineral particles, without intimate mixing, in the underlying mineral horizons. Classification One of the first soil classification systems was developed by Russian scientist Vasily Dokuchaev around 1880. It was modified a number of times by American and European researchers and was developed into the system commonly used until the 1960s. It was based on the idea that soils have a particular morphology based on the materials and factors that form them. In the 1960s, a different classification system began to emerge which focused on soil morphology instead of parental materials and soil-forming factors. Since then, it has undergone further modifications. The World Reference Base for Soil Resources aims to establish an international reference base for soil classification. Uses Soil is used in agriculture, where it serves as the anchor and primary nutrient base for plants. The types of soil and available moisture determine the species of plants that can be cultivated. Agricultural soil science was the primeval domain of soil knowledge, long time before the advent of pedology in the 19th century. However, as demonstrated by aeroponics, aquaponics and hydroponics, soil material is not an absolute essential for agriculture, and soilless cropping systems have been claimed as the future of agriculture for an endless growing mankind. Soil material is also a critical component in mining, construction and landscape development industries. Soil serves as a foundation for most construction projects. The movement of massive volumes of soil can be involved in surface mining, road building and dam construction. Earth sheltering is the architectural practice of using soil for external thermal mass against building walls. Many building materials are soil based. Loss of soil through urbanization is growing at a high rate in many areas and can be critical for the maintenance of subsistence agriculture. Soil resources are critical to the environment, as well as to food and fibre production, producing 98.8% of food consumed by humans. Soil provides minerals and water to plants according to several processes involved in plant nutrition. Soil absorbs rainwater and releases it later, thus preventing floods and drought, flood regulation being one of the major ecosystem services provided by soil. Soil cleans water as it percolates through it. Soil is the habitat for many organisms: the major part of known and unknown biodiversity is in the soil, in the form of earthworms, woodlice, millipedes, centipedes, snails, slugs, mites, springtails, enchytraeids, nematodes, protists), bacteria, archaea, fungi and algae; and most organisms living above ground have part of them (plants) or spend part of their life cycle (insects) below-ground. Above-ground and below-ground biodiversities are tightly interconnected, making soil protection of paramount importance for any restoration or conservation plan. The biological component of soil is an extremely important carbon sink since about 57% of the biotic content is carbon. Even in deserts, cyanobacteria, lichens and mosses form biological soil crusts which capture and sequester a significant amount of carbon by photosynthesis. Poor farming and grazing methods have degraded soils and released much of this sequestered carbon to the atmosphere. Restoring the world's soils could offset the effect of increases in greenhouse gas emissions and slow global warming, while improving crop yields and reducing water needs. Waste management often has a soil component. Septic drain fields treat septic tank effluent using aerobic soil processes. Land application of waste water relies on soil biology to aerobically treat BOD. Alternatively, landfills use soil for daily cover, isolating waste deposits from the atmosphere and preventing unpleasant smells. Composting is now widely used to treat aerobically solid domestic waste and dried effluents of settling basins. Although compost is not soil, biological processes taking place during composting are similar to those occurring during decomposition and humification of soil organic matter. Organic soils, especially peat, serve as a significant fuel and horticultural resource. Peat soils are also commonly used for the sake of agriculture in Nordic countries, because peatland sites, when drained, provide fertile soils for food production. However, wide areas of peat production, such as rain-fed sphagnum bogs, also called blanket bogs or raised bogs, are now protected because of their patrimonial interest. As an example, Flow Country, covering 4,000 square kilometres of rolling expanse of blanket bogs in Scotland, is now candidate for being included in the World Heritage List. Under present-day global warming peat soils are thought to be involved in a self-reinforcing (positive feedback) process of increased emission of greenhouse gases (methane and carbon dioxide) and increased temperature, a contention which is still under debate when replaced at field scale and including stimulated plant growth. Geophagy is the practice of eating soil-like substances. Both animals and humans occasionally consume soil for medicinal, recreational, or religious purposes. It has been shown that some monkeys consume soil, together with their preferred food (tree foliage and fruits), in order to alleviate tannin toxicity. Soils filter and purify water and affect its chemistry. Rain water and pooled water from ponds, lakes and rivers percolate through the soil horizons and the upper rock strata, thus becoming groundwater. Pests (viruses) and pollutants, such as persistent organic pollutants (chlorinated pesticides, polychlorinated biphenyls), oils (hydrocarbons), heavy metals (lead, zinc, cadmium), and excess nutrients (nitrates, sulfates, phosphates) are filtered out by the soil. Soil organisms metabolise them or immobilise them in their biomass and necromass, thereby incorporating them into stable humus. The physical integrity of soil is also a prerequisite for avoiding landslides in rugged landscapes. Degradation Land degradation is a human-induced or natural process which impairs the capacity of land to function. Soil degradation involves acidification, contamination, desertification, erosion or salination. Acidification Soil acidification is beneficial in the case of alkaline soils, but it degrades land when it lowers crop productivity, soil biological activity and increases soil vulnerability to contamination and erosion. Soils are initially acid and remain such when their parent materials are low in basic cations (calcium, magnesium, potassium and sodium). On parent materials richer in weatherable minerals acidification occurs when basic cations are leached from the soil profile by rainfall or exported by the harvesting of forest or agricultural crops. Soil acidification is accelerated by the use of acid-forming nitrogenous fertilizers and by the effects of acid precipitation. Deforestation is another cause of soil acidification, mediated by increased leaching of soil nutrients in the absence of tree canopies. Contamination Soil contamination at low levels is often within a soil's capacity to treat and assimilate waste material. Soil biota can treat waste by transforming it, mainly through microbial enzymatic activity. Soil organic matter and soil minerals can adsorb the waste material and decrease its toxicity, although when in colloidal form they may transport the adsorbed contaminants to subsurface environments. Many waste treatment processes rely on this natural bioremediation capacity. Exceeding treatment capacity can damage soil biota and limit soil function. Derelict soils occur where industrial contamination or other development activity damages the soil to such a degree that the land cannot be used safely or productively. Remediation of derelict soil uses principles of geology, physics, chemistry and biology to degrade, attenuate, isolate or remove soil contaminants to restore soil functions and values. Techniques include leaching, air sparging, soil conditioners, phytoremediation, bioremediation and Monitored Natural Attenuation. An example of diffuse pollution with contaminants is copper accumulation in vineyards and orchards to which fungicides are repeatedly applied, even in organic farming. Microfibres from synthetic textiles are another type of plastic soil contamination, 100% of agricultural soil samples from southwestern China contained plastic particles, 92% of which were microfibres. Sources of microfibres likely included string or twine, as well as irrigation water in which clothes had been washed. The application of biosolids from sewage sludge and compost can introduce microplastics to soils. This adds to the burden of microplastics from other sources (e.g. the atmosphere). Approximately half the sewage sludge in Europe and North America is applied to agricultural land. In Europe it has been estimated that for every million inhabitants 113 to 770 tonnes of microplastics are added to agricultural soils each year. Desertification Desertification, an environmental process of ecosystem degradation in arid and semi-arid regions, is often caused by badly adapted human activities such as overgrazing or excess harvesting of firewood. It is a common misconception that drought causes desertification. Droughts are common in arid and semiarid lands. Well-managed lands can recover from drought when the rains return. Soil management tools include maintaining soil nutrient and organic matter levels, reduced tillage and increased cover. These practices help to control erosion and maintain productivity during periods when moisture is available. Continued land abuse during droughts, however, increases land degradation. Increased population and livestock pressure on marginal lands accelerates desertification. It is now questioned whether present-day climate warming will favour or disfavour desertification, with contradictory reports about predicted rainfall trends associated with increased temperature, and strong discrepancies among regions, even in the same country. Erosion Erosion of soil is caused by water, wind, ice, and movement in response to gravity. More than one kind of erosion can occur simultaneously. Erosion is distinguished from weathering, since erosion also transports eroded soil away from its place of origin (soil in transit may be described as sediment). Erosion is an intrinsic natural process, but in many places it is greatly increased by human activity, especially unsuitable land use practices. These include agricultural activities which leave the soil bare during times of heavy rain or strong winds, overgrazing, deforestation, and improper construction activity. Improved management can limit erosion. Soil conservation techniques which are employed include changes of land use (such as replacing erosion-prone crops with grass or other soil-binding plants), changes to the timing or type of agricultural operations, terrace building, use of erosion-suppressing cover materials (including cover crops and other plants), limiting disturbance during construction, and avoiding construction during erosion-prone periods and in erosion-prone places such as steep slopes. Historically, one of the best examples of large-scale soil erosion due to unsuitable land-use practices is wind erosion (the so-called dust bowl) which ruined American and Canadian prairies during the 1930s, when immigrant farmers, encouraged by the federal government of both countries, settled and converted the original shortgrass prairie to agricultural crops and cattle ranching. A serious and long-running water erosion problem occurs in China, on the middle reaches of the Yellow River and the upper reaches of the Yangtze River. From the Yellow River, over 1.6 billion tons of sediment flow each year into the ocean. The sediment originates primarily from water erosion (gully erosion) in the Loess Plateau region of northwest China. Soil piping is a particular form of soil erosion that occurs below the soil surface. It causes levee and dam failure, as well as sink hole formation. Turbulent flow removes soil starting at the mouth of the seep flow and the subsoil erosion advances up-gradient. The term sand boil is used to describe the appearance of the discharging end of an active soil pipe. Salination Soil salination is the accumulation of free salts to such an extent that it leads to degradation of the agricultural value of soils and vegetation. Consequences include corrosion damage, reduced plant growth, erosion due to loss of plant cover and soil structure, and water quality problems due to sedimentation. Salination occurs due to a combination of natural and human-caused processes. Arid conditions favour salt accumulation. This is especially apparent when soil parent material is saline. Irrigation of arid lands is especially problematic. All irrigation water has some level of salinity. Irrigation, especially when it involves leakage from canals and overirrigation in the field, often raises the underlying water table. Rapid salination occurs when the land surface is within the capillary fringe of saline groundwater. Soil salinity control involves watertable control and flushing with higher levels of applied water in combination with tile drainage or another form of subsurface drainage. Reclamation Soils which contain high levels of particular clays with high swelling properties, such as smectites, are often very fertile. For example, the smectite-rich paddy soils of Thailand's Central Plains are among the most productive in the world. However, the overuse of mineral nitrogen fertilizers and pesticides in irrigated intensive rice production has endangered these soils, forcing farmers to implement integrated practices based on Cost Reduction Operating Principles. Many farmers in tropical areas, however, struggle to retain organic matter and clay in the soils they work. In recent years, for example, productivity has declined and soil erosion has increased in the low-clay soils of northern Thailand, following the abandonment of shifting cultivation for a more permanent land use. Farmers initially responded by adding organic matter and clay from termite mound material, but this was unsustainable in the long-term because of rarefaction of termite mounds. Scientists experimented with adding bentonite, one of the smectite family of clays, to the soil. In field trials, conducted by scientists from the International Water Management Institute (IWMI) in cooperation with Khon Kaen University and local farmers, this had the effect of helping retain water and nutrients. Supplementing the farmer's usual practice with a single application of of bentonite resulted in an average yield increase of 73%. Other studies showed that applying bentonite to degraded sandy soils reduced the risk of crop failure during drought years. In 2008, three years after the initial trials, IWMI scientists conducted a survey among 250 farmers in northeast Thailand, half of whom had applied bentonite to their fields. The average improvement for those using the clay addition was 18% higher than for non-clay users. Using the clay had enabled some farmers to switch to growing vegetables, which need more fertile soil. This helped to increase their income. The researchers estimated that 200 farmers in northeast Thailand and 400 in Cambodia had adopted the use of clays, and that a further 20,000 farmers were introduced to the new technique. If the soil is too high in clay or salts (e.g. saline sodic soil), adding gypsum, washed river sand and organic matter (e.g.municipal solid waste) will balance the composition. Adding organic matter, like ramial chipped wood or compost, to soil which is depleted in nutrients and too high in sand will boost its quality and improve production. Special mention must be made of the use of charcoal, and more generally biochar to improve nutrient-poor tropical soils, a process based on the higher fertility of anthropogenic pre-Columbian Amazonian Dark Earths, also called Terra Preta de Índio, due to interesting physical and chemical properties of soil black carbon as a source of stable humus. However, the uncontrolled application of charred waste products of all kinds may endanger soil life and human health. History of studies and research The history of the study of soil is intimately tied to humans' urgent need to provide food for themselves and forage for their animals. Throughout history, civilizations have prospered or declined as a function of the availability and productivity of their soils. Studies of soil fertility The Greek historian Xenophon (450–355 BCE) was the first to expound upon the merits of green-manuring crops: 'But then whatever weeds are upon the ground, being turned into earth, enrich the soil as much as dung.' Columella's Of husbandry, , advocated the use of lime and that clover and alfalfa (green manure) should be turned under, and was used by 15 generations (450 years) under the Roman Empire until its collapse. From the fall of Rome to the French Revolution, knowledge of soil and agriculture was passed on from parent to child and as a result, crop yields were low. During the European Middle Ages, Yahya Ibn al-'Awwam's handbook, with its emphasis on irrigation, guided the people of North Africa, Spain and the Middle East; a translation of this work was finally carried to the southwest of the United States when under Spanish influence. Olivier de Serres, considered the father of French agronomy, was the first to suggest the abandonment of fallowing and its replacement by hay meadows within crop rotations. He also highlighted the importance of soil (the French terroir) in the management of vineyards. His famous book contributed to the rise of modern, sustainable agriculture and to the collapse of old agricultural practices such as soil amendment for crops by the lifting of forest litter and assarting, which ruined the soils of western Europe during the Middle Ages and even later on according to regions. Experiments into what made plants grow first led to the idea that the ash left behind when plant matter was burned was the essential element but overlooked the role of nitrogen, which is not left on the ground after combustion, a belief which prevailed until the 19th century. In about 1635, the Flemish chemist Jan Baptist van Helmont thought he had proved water to be the essential element from his famous five years' experiment with a willow tree grown with only the addition of rainwater. His conclusion came from the fact that the increase in the plant's weight had apparently been produced only by the addition of water, with no reduction in the soil's weight. John Woodward ( 1728) experimented with various types of water ranging from clean to muddy and found muddy water the best, and so he concluded that earthy matter was the essential element. Others concluded it was humus in the soil that passed some essence to the growing plant. Still others held that the vital growth principal was something passed from dead plants or animals to the new plants. At the start of the 18th century, Jethro Tull demonstrated that it was beneficial to cultivate (stir) the soil, but his opinion that the stirring made the fine parts of soil available for plant absorption was erroneous. As chemistry developed, it was applied to the investigation of soil fertility. The French chemist Antoine Lavoisier showed in about 1778 that plants and animals must combust oxygen internally to live. He was able to deduce that most of the weight of van Helmont's willow tree derived from air. It was the French agriculturalist Jean-Baptiste Boussingault who by means of experimentation obtained evidence showing that the main sources of carbon, hydrogen and oxygen for plants were air and water, while nitrogen was taken from soil. Justus von Liebig in his book Organic chemistry in its applications to agriculture and physiology (published 1840), asserted that the chemicals in plants must have come from the soil and air and that to maintain soil fertility, the used minerals must be replaced. Liebig nevertheless believed the nitrogen was supplied from the air. The enrichment of soil with guano by the Incas was rediscovered in 1802, by Alexander von Humboldt. This led to its mining and that of Chilean nitrate and to its application to soil in the United States and Europe after 1840. The work of Liebig was a revolution for agriculture, and so other investigators started experimentation based on it. In England John Bennet Lawes and Joseph Henry Gilbert worked in the Rothamsted Experimental Station, founded by the former, and that plants took nitrogen from the soil, and that salts needed to be in an available state to be absorbed by plants. Their investigations also produced the superphosphate, consisting in the acid treatment of phosphate rock. This led to the invention and use of salts of potassium (K) and nitrogen (N) as fertilizers. Ammonia generated by the production of coke was recovered and used as fertiliser. Finally, the chemical basis of nutrients delivered to the soil in manure was understood and in the mid-19th century chemical fertilisers were applied. However, the dynamic interaction of soil and its life forms was still not understood. In 1856, J. Thomas Way discovered that ammonia contained in fertilisers was transformed into nitrates, and twenty years later Robert Warington proved that this transformation was done by living organisms. In 1890 Sergei Winogradsky announced he had found the bacteria responsible for this transformation. It was known that certain legumes could take up nitrogen from the air and fix it to the soil but it took the development of bacteriology towards the end of the 19th century to lead to an understanding of the role played in nitrogen fixation by bacteria. The symbiosis of bacteria and leguminous roots, and the fixation of nitrogen by the bacteria, were simultaneously discovered by the German agronomist Hermann Hellriegel and the Dutch microbiologist Martinus Beijerinck. Crop rotation, mechanisation, chemical and natural fertilisers led to a doubling of wheat yields in western Europe between 1800 and 1900. Studies of soil formation Scientists who studied soil in connection with agricultural practices considered it mainly a static substrate. However, the soil is the result of evolution from more ancient geological materials under the action of biotic and abiotic processes. After studies of soil improvement commenced, other researchers began to study soil genesis and, as a result, soil types and classifications. In 1860, while in Mississippi, Eugene W. Hilgard (1833–1916) studied the relationship between rock material, climate, vegetation, and the type of soils that were developed. He realised that the soils were dynamic and considered the classification of soil types. (	Physical sciences	Earth science	null
37744	https://en.wikipedia.org/wiki/M1%20Abrams	M1 Abrams	The M1 Abrams () is a third-generation American main battle tank designed by Chrysler Defense (now General Dynamics Land Systems) and named for General Creighton Abrams. Conceived for modern armored ground warfare, it is one of the heaviest tanks in service at nearly . It introduced several modern technologies to the United States armored forces, including a multifuel turbine engine, sophisticated Chobham composite armor, a computer fire control system, separate ammunition storage in a blowout compartment, and NBC protection for crew safety. Initial models of the M1 were armed with a 105 mm M68 gun, while later variants feature a license-produced Rheinmetall 120 mm L/44 designated M256. The M1 Abrams was developed from the failed joint American-West German MBT-70 project that intended to replace the dated M60 tank. There are three main operational Abrams versions: the M1, M1A1, and M1A2, with each new iteration seeing improvements in armament, protection, and electronics. The Abrams was to be replaced in U.S. Army service by the XM1202 Mounted Combat System, but because that project was canceled, the Army has opted to continue maintaining and operating the M1 series for the foreseeable future by upgrading optics, armor, and firepower. The M1 Abrams entered service in 1980 and serves as the main battle tank of the United States Army and formerly of the U.S. Marine Corps (USMC) until the decommissioning of all USMC tank battalions in 2021. The export modification is used by the armed forces of Egypt, Kuwait, Saudi Arabia, Australia, Poland and Iraq. The Abrams was first used in combat by the U.S. in the Gulf War. It was later deployed by the U.S. in the War in Afghanistan and the Iraq War, as well as by Iraq in the war against the Islamic State, Saudi Arabia in the Yemeni Civil War, and Ukraine during the Russian invasion of Ukraine. History Previous developments In 1963, the U.S. Army and the West German Bundeswehr began collaborating on a main battle tank (MBT) design that both nations would use, improving interoperability between the two NATO partners. The MBT-70, or Kampfpanzer 70 as it was known in Germany, incorporated many new unconventional technologies across the board. Conventional tanks of the time had a crew of four, with the driver located in the hull. In the MBT-70, the loader crewmember would be replaced by a mechanical autoloader and the driver would be located inside the NBC-protected turret with the other two crewmembers. Like the M60A2 MBT and M551 Sheridan light tank then under development, the MBT-70 was armed with a 152 mm gun-launcher that, in addition to firing conventional ammunition, would also fire the Shillelagh missile. A hydropneumatic suspension provided improved cross-country ride quality and also allowed the entire tank to be raised or lowered by the driver. The United States team was led by General Motors while the German team consisted of a consortium of firms. The collaboration between the two teams was rocky from the start, with many cultural differences and disagreements about the design hampering progress. Germany favored a tank optimized for the terrain of central Europe while the U.S. attached importance to operating anywhere in the world. The Germans had reservations about the Shillelagh missile and developed a 120 mm high-velocity gun as an alternative. Perhaps the most contentious disagreement, never fully resolved, concerned the measurement system to be used in drafting. Germany became concerned with the excessive weight of the tank. In light of growing costs, delays and overall uncertainty as to the soundness of the tank design, the United States and Germany ended their MBT-70 partnership in 1970. The U.S. Army began work on an austere version of the MBT-70, named XM803. Systems were simplified or eliminated altogether and the unreliable autoloader was improved. These changes were ultimately insufficient to allay concerns about the tank's cost. Congress canceled the XM803 in December 1971 but permitted the Army to reallocate remaining funds to develop a new main battle tank. Starting afresh The Army began the XM815 project in January 1972. The Main Battle Tank Task Force (MBTTF) was established under Major General William Desobry. The task force prepared design studies with the technical support of Tank-automotive and Armaments Command (TACOM). TACOM began examining specific goals. To this end, a new design basis emerged in February 1973. It had to defeat any hit from a Soviet gun within and 30 degrees to either side. The tank would be armed with the 105 mm M68 gun, a licensed version of the Royal Ordnance L7, and a 20 mm version of the M242 Bushmaster. The Army later deleted the latter from the design, seeing it as superfluous. In spring 1972, Desobry was briefed by the British on their own newly developed "Burlington" armor from the British Army's labs. The armor performed exceptionally against shaped charges such as HEAT rounds. In September, Desobry convinced the Army to incorporate the new armor. To take full advantage of Burlington, also known as Chobham, the new tank would have to have armor around two feet thick (for comparison, the armor on the M60 is around four inches thick). General Creighton Abrams set the weight of the new tank at . The original goal of keeping weight under was abandoned. At the time, the Pentagon's procurement system was beset with problems being caused by the desire to have the best possible design. This often resulted in programs being canceled due to cost overruns, leaving the forces with outdated systems, as was the case with the MBT-70. There was a strong movement within the Army to get a new design within budget to prevent the MBT-70 experience from repeating itself. For the new design, the Army set the design-to-unit cost at no more than $507,790 (). The Pentagon's approach to control of research and development was modified with the XM1. Previous acquisition strategy called for a significant amount of the design work to be done by the government. Under the new framework, contractors would competitively bid their own designs rather than compete solely for the right to manufacture the end product. In January 1973, the U.S. Army issued the XM1 (as the XM815 had been renamed in November 1972) request for proposals. In May 1973, Chrysler Defense and General Motors submitted proposals. Both were armed with the 105 mm M68 gun, the licensed L7, and the 20 mm Bushmaster. Chrysler chose a 1,500 hp Lycoming AGT1500 gas turbine engine. GM's model was powered by a 1,500 hp diesel engine similar to that used on the American MBT-70 and XM803. Prototypes Prototypes were delivered in 1976 by Chrysler and GM armed with the M68E1 105 mm gun. They entered head-to-head testing at Aberdeen Proving Ground. The testing showed that the GM design was generally superior to Chrysler's, offering better armor protection, and better fire control and turret stabilization systems. During testing, the power packs of both designs proved to have issues. The Chrysler gas turbine engine had extensive heat recovery systems in an attempt to improve its fuel efficiency to something similar to a traditional internal combustion engine. This goal was not achieved: the engine consumed much more fuel than expected, burning . The GM design used a new variable-compression diesel design. By spring 1976, the decision to choose the GM design was largely complete. In addition to offering better overall performance, there were concerns about Chrysler's engine both from a reliability and fuel consumption standpoint. The GM program was also slightly cheaper overall at $208 million compared to $221 million for Chrysler. In July 1976, the Army prepared to inform Congress of the decision to move ahead with the GM design. All that was required was the final sign-off by the U.S Secretary of Defense, Donald Rumsfeld. Back to the drawing board On 20 July 1976, United States Secretary of the Army Martin Hoffmann and a group of generals visited Deputy Defense Secretary Bill Clements and Director of Defense Research and Engineering Malcolm Currie on their decision. They were surprised when Clements and Currie criticized their decision and demanded that the new tank have a turbine. Defense Secretary Rumsfeld heard arguments from both parties in the afternoon. The Army team spent the night writing briefs and presented them to Rumsfeld the next morning, who then announced a four-month delay. Within days, GM was asked to present a new design with a turbine engine. According to Assistant Secretary for Research and Development Ed Miller, "It became increasingly clear that the only solution which would be acceptable to Clements and Currie was the turbine... It was a political decision that was reached, and for all intents and purposes that decision gave the award to Chrysler since they were the only contractor with a gas turbine." In the meantime, in September 1976 three West German Leopard 2AV prototypes were belatedly sent to Aberdeen for comparison testing. Germany had signed a somewhat vague memorandum of understanding in 1974 committing both parties toward commonality in tank parts. Germany had assumed that its tank would be evaluated against the GM and Chrysler's prototypes and that the best tank would be chosen for production. This misunderstanding arose from the fact that in public statements both countries had overrepresented the MOU as an agreement that Germany and the U.S. would select a common MBT. In reality, the U.S. Army was unwilling to choose a foreign tank unless it was obviously superior in design and cost. In any case, in evaluations the Leopard 2AV was found to meet U.S. requirements but was thought to cost more. The U.S. Army announced in January 1977 that Germany had withdrawn the tank from consideration. Chrysler is chosen Having narrowly averted losing the contract, Chrysler set about improving the design. Expensive components were replaced with less expensive ones. Chrysler's team also negotiated lower costs from their subcontractors. The price of the redesigned tank's turret especially was decreased, but other improvements came from unexpected places, such as a $600 hydraulic oil reservoir replaced with a $25 one. Chrysler also submitted a version with a Teledyne AVCR-1360 diesel engine. Chrysler's new bid came to $196 million, down from $221 million in the original proposal. GM's proposal replaced the diesel engine with an AGT1500 turbine and integrated a turret capable of mounting either the 105 mm or 120 mm gun. Cost growth pushed the tank bid to $232 million from $208 million. Although the GM team had successfully integrated the turbine, Baer was more impressed by the cost savings introduced by the Chrysler team's redesign. On 12 November 1976, the Defense Department awarded the $4.9 billion development contract to Chrysler. The turbine engine and cost do not appear to be the only reason for the selection of Chrysler. Chrysler was the only company that appeared to be seriously interested in tank development; the M60 had been lucrative for the company. In contrast, GM made only about 1% of its income from military sales, compared to 5% for Chrysler, and only submitted their bid after a "special plea" from the Pentagon. Eleven XM1 preproduction models were manufactured between February and July 1978 at Detroit Arsenal Tank Plant. Quality problems with the engine quickly became apparent in testing. The first preproduction units that arrived at Aberdeen Proving Ground in March 1978 had serious problems. The tank accumulated mud and dirt under the hull which led to thrown tracks. Chrysler installed a scraper to prevent the build-up of dirt. This did not solve the issue entirely. It was determined months later that a gauge used to tension tracks was miscalibrated. This caused the tracks to be fitted too loosely. Another problem was the ingestion of debris by the engine. The problem was determined to be caused by poorly fitting air filters. At Fort Bliss, several tanks experienced transmission issues. It was determined that the tankers at Fort Bliss had discovered that they could throw the vehicle directly from acceleration into reverse, a tactically advantageous maneuver called the "bow tie". Chrysler resolved this by installing a device that prevented this. The problems found during testing were easily surmounted. Critics of the M1 program emerged in the early 1980s, particularly the newly formed Project on Military Procurement (PMP) (later renamed the Project on Government Oversight). PMP took issue with the tank's vulnerability, high price, reliance on flammable hydraulics, and high fuel consumption. American tank historian Steven J. Zaloga characterized American press criticism of the M1 during this time as "ill-founded". Zaloga wrote the issues uncovered by the tank trials were "not particularly serious". PMP's criticism failed to generate any serious opposition to the program, which maintained strong support from Congress and the Pentagon. Responding to some of the alleged issues with the tank in King of the Killing Zone (1989), journalist Orr Kelly wrote that "The truth is close to the opposite." Kelly said the program "ranks as one of the Army's best managed", producing a tank in "a remarkably short time" while avoiding "gold-plating" and utilizing effective competition. Production starts Low rate initial production (LRIP) of the vehicle was approved in May 1979. In February 1982, General Dynamics Land Systems Division (GDLS) purchased Chrysler Defense, after Chrysler built over 1,000 M1s. A total of 3,273 M1 Abrams tanks were produced during 1979–1985 and first entered U.S. Army service in 1980. Production at the government-owned, GDLS-operated Lima Army Tank Plant in Lima, Ohio, was joined by vehicles built at the Detroit Arsenal Tank Plant (DATP) in Warren, Michigan from 1982 to 1991 (DATP also produced the 11 preproduction models in 1978.). The U.S. Army Laboratory Command (LABCOM), under the supervision of the United States Army Research Laboratory (ARL), was also heavily involved with designing the tank with M1A1 armor resistant shells, M829A2 armor-penetrating rounds, and improved weapon range. The M1 was armed with the license-built M68A1 version of the 105 mm Royal Ordnance L7 gun. The tank featured the first-of-its-kind Chobham armor. The M1 Abrams was the first to use this advanced armor. It consisted of an arrangement of metal and ceramic plates. An improved model called the IPM1 was produced briefly in 1984 and contained upgrades to armor and other small improvements. 120 mm gun M1A1 A number of considerations had led the service and its contractors to favor the Army's standard M68 105 mm gun over Germany's 120 mm Rheinmetall Rh-120 smoothbore gun for the XM1. To begin with, the 105 mm gun was "the smallest, lightest, and least costly gun adequate for the job." Indeed, new kinetic energy ammunition for the weapon then under development by the Army promised to extend the gun's usefulness well into the future. And because the Army's other tanks, the M60 and the upgraded M48, as well as the tanks of virtually every other NATO nation, used the 105 mm gun, mounting that gun on the XM1 promised to increase standardization within the alliance. Moreover, the continuing development of the new ammunition for the XM1 automatically upgraded every other gun in NATO. For all of these reasons, the XM1's development proceeded "on the assumption that the 105 mm gun would probably be the eventual main armament." The tripartite British—American—German gun trials of 1975 produced a general agreement in the U.S. Defense Department that at some future point, a 120 mm gun of some design would be added to the XM1. Apparently anticipating this, Chrysler and GM had both made changes to their tanks during development to make them compatible with a variety of main guns. In January 1978, the Secretary of the Army announced that the Rheinmetall 120 mm gun would be mounted on future production versions of the XM1. This decision established the requirement for a separate program for the M1E1 (with 120 mm gun) so that the XM1 program could continue unimpeded. About 5,000 M1A1 Abrams tanks were produced from 1986 to 1992 and featured the M256 120 mm smoothbore cannon, improved armor, consisting of depleted uranium and other classified materials, and a CBRN protection system. Production of M1 and M1A1 tanks totaled some 9,000 tanks at a cost of approximately $4.3 million per unit. In 1990, a Project On Government Oversight report criticized the M1's high costs and low fuel efficiency in comparison with other tanks of similar power and effectiveness such as the Leopard 2. As the Abrams entered service, they operated alongside M60A3 within the U.S. military and with other NATO tanks in various Cold War exercises which usually took place in Western Europe, especially West Germany. The exercises were aimed at countering Soviet forces. Adaptations before the Gulf War (Operations Desert Shield and Desert Storm) gave the vehicle better firepower and Nuclear, Biological and Chemical (NBC) protection. Gulf War The Abrams remained untested in combat until the Gulf War in 1991, during Operation Desert Storm. The first Abrams tanks to arrive in Saudi Arabia in August 1990 in the buildup to the war were M1 and IPM1 tanks with 105 mm guns. All but two battalions of 105 mm gun Abrams tanks were replaced by M1A1 tanks prior to the American invasion in January 1991. The U.S. Army deployed a total of 1,956 M1A1s (733 M1A1, 1,233 M1A1HA) to Saudi Arabia to participate in the liberation of Kuwait. The U.S. Marine Corps deployed 353 tanks, of which 277 were M60s and 76 were M1A1 (60 M1A1HA and 16 M1A1 Common). The M1A1 Common variant included adaptations for deep wading and improvements to increase commonality with the Army's Abrams. The 2nd Tank Battalion was equipped with M1A1HA Abrams borrowed from the Army. The M1A1 was superior to Iraq's Soviet-era T-54/T-55 and T-62 tanks, as well as T-72 versions imported from the Soviet Union and Poland. Polish officials stated that no license-produced T-72 (nicknamed Lion of Babylon) tanks were finished before destruction of the Iraqi Taji tank plant in 1991. Iraq's T-72s, like most Soviet export designs, lacked night-vision systems and then-modern rangefinders, though they did have some night-fighting tanks with older active infrared systems or floodlights. Very few M1 tanks were hit by enemy fire and none were destroyed as a direct result of enemy fire, none of which resulted in any fatalities. Three Abrams were left behind the enemy lines after a swift attack on Talil airfield, south of Nasiriyah, on February 27. One of them was hit by enemy fire, while the other two became embedded in mud. The tanks were destroyed by U.S. forces to prevent any trophy-claim by the Iraqi Army. A total of 23 M1A1s were damaged or destroyed during the war. Of the nine Abrams tanks destroyed, seven were destroyed by friendly fire and two intentionally destroyed to prevent capture by the Iraqi Army. No M1s were lost to enemy tank fire. Some others took minor combat damage, with little effect on their operational readiness. The M1A1 could kill other tanks at ranges in excess of . This range was crucial in combat against previous generation tanks of Soviet design in Desert Storm, as the effective range of the main gun in the Iraqi tanks was less than . This meant Abrams tanks could hit Iraqi tanks before the enemy got in range—a decisive advantage in this kind of combat. In friendly fire incidents, the front armor and fore side turret armor survived direct APFSDS hits from other M1A1s. This was not the case for the side armor of the hull and the rear armor of the turret, as both areas were penetrated on at least two occasions by unintentional strikes by depleted uranium ammunition during the Battle of Norfolk. Waco siege During the Waco siege in 1993, two M1A1 Abrams tanks were borrowed from the military and deployed by the FBI against the Branch Davidians. Upgrades The M1A2 was a further improvement of the M1A1, with a commander's independent thermal viewer, weapon station, position navigation equipment, and a full set of controls and displays linked by a digital data bus. These upgrades also provided the M1A2 with an improved fire control system. The M1A2 System Enhancement Package (SEP) added digital maps, Force XXI Battle Command Brigade and Below (FBCB2) Linux communications system capabilities for commanders, and an improved cooling system to compensate for heat generated by the additional computer systems. The M1A2 SEP also serves as the basis for the M104 Wolverine heavy assault bridge. The M1A2 SEPv2 (version 2) added Common Remotely Operated Weapon Station (CROWS or CROWS II) support, color displays, better interfaces, a new operating system, better front and side armor, and an upgraded transmission for better durability. Further upgrades included depleted uranium armor for all variants, a system overhaul that returns all A1s to like-new condition (M1A1 AIM), a digital enhancement package for the A1 (M1A1D), and a commonality program to standardize parts between the U.S. Army and the Marine Corps (M1A1HC). Improvements to survivability, lethality, and protection have been sought since 2014. Iraq War Further combat was seen during 2003 when U.S. forces invaded Iraq and deposed Iraqi President Saddam Hussein in the Iraq War's Operation Iraqi Freedom. One achievement of the M1A1s was the destruction of seven T-72s in a point-blank skirmish (less than ) near Mahmoudiyah, about south of Baghdad, with no U.S. losses. This was in the face of inadequately trained Iraqi tank crews, most of whom had not fired live ammunition in the previous year due to the sanctions then in operation and made no hits at point-blank range. Following lessons learned in Desert Storm, the Abrams and many other U.S. combat vehicles used in the conflict were fitted with Combat Identification Panels to reduce friendly fire incidents. Several Abrams tanks that were irrecoverable due to loss of mobility or other circumstances were destroyed by friendly forces, usually by other Abrams tanks, to prevent their capture. Some Abrams tanks were disabled by Iraqi infantrymen in ambushes during the invasion. Some troops employed short-range anti-tank rockets and fired at the tracks, rear and top. Other tanks were put out of action by engine fires when flammable fuel stored externally in turret racks was hit by small arms fire and spilled into the engine compartment. By March 2005, approximately 80 Abrams tanks had been forced out of action by enemy attacks; 63 were shipped back to the U.S. for repairs, while 17 were damaged beyond repair with 3 of them at the beginning of 2003. Vulnerabilities exposed during urban combat in the Iraq War were addressed with the Tank Urban Survival Kit (TUSK) modifications, including armor upgrades and a gun shield, issued to some M1 Abrams tanks. It added protection in the rear and side of the tank and improved fighting ability and survival ability in urban environments. By December 2006 more than 530 Abrams tanks had been shipped back to the U.S. for repairs. In May 2008, it was reported that a U.S. M1 tank had also been damaged in Iraq by insurgent fire of a Soviet-made RPG-29 "Vampir", which uses a tandem-charge HEAT warhead to penetrate explosive reactive armor (ERA) as well as composite armor behind it. The U.S. considered the RPG-29 a high threat to armor and refused to allow the newly formed Iraqi Army to buy it, fearing that it would fall into the insurgents' hands. Iraqi Army service Between 2010 and 2012 the U.S. supplied 140 refurbished M1A1 Abrams tanks to Iraq. In mid-2014, they saw action when the Islamic State of Iraq and the Levant (ISIL or Islamic State) launched the June 2014 Northern Iraq offensive. During three months, about one-third of the Iraqi Army's M1 tanks had been damaged or destroyed by ISIL and some were captured by opposing forces. By December 2014, the Iraqi Army only had about 40 operational Abrams left. That month, the U.S. Department of State approved the sale of another 175 Abrams to Iraq. Iranian-backed Iraqi Shiite Kata'ib Hezbollah (Hezbollah Brigades) were reported to operate M1 Abrams, and released publicity showing the tanks being transported by trucks to take part in the Battle of Mosul. It is not known whether the tanks were captured from ISIL, seized from Iraq's military, or handed over. One Iraqi-operated Abrams has been nicknamed "The Beast" after it became the lone working tank when taking back the town of Hit in April 2016, destroying enemy fighting positions and IED emplacements. In October 2017, Abrams were used by the Iraqi security forces and the Popular Mobilization Forces (also called Al-Hashd al-Shaabi) in assaults against the Kurdistan Regional Government Peshmerga in the town of Altun Kupri (also called Prde). It was claimed by Kurdish commanders that at least one Abrams was destroyed by the Peshmerga. War in Afghanistan Canada and Denmark deployed Leopard 1 and 2 MBTs that were specially modified to operate in the relatively flat and arid conditions of southwestern Afghanistan. In late 2010, at the request of Regional Command Southwest, the U.S. Marine Corps deployed a small detachment of 14 M1A1 Abrams tanks from Delta Company, 1st Tank Battalion, 1st Marine Division (Forward), to southern Afghanistan in support of operations in Helmand and Kandahar provinces. 2015 Yemen Civil War Saudi Abrams tanks saw service in the 2015 Yemeni Civil War, where M1A2s were used against Houthi rebels. In August 2016, the U.S. approved a deal to sell up to 153 more Abrams tanks to Saudi Arabia, including 20 "battle damage replacements", suggesting that some Saudi Arabian Abrams had been destroyed or severely damaged in combat in Yemen. Russo-Ukrainian War Russian invasion of Ukraine In January 2023, U.S. President Joe Biden said that the United States would send 31 M1 Abrams tanks to Ukraine. The plan to transfer the tanks to Ukraine was approved as part of a larger aid package. Pentagon spokesperson Sabrina Singh specified that the tanks would be the M1A2 variant; however, because they were not available in excess in U.S. stocks, they would be purchased through Ukraine Security Assistance Initiative (USAI) and could take up to two years to manufacture and deliver. In March 2023 the Pentagon announced that, in order to expedite delivery, older M1A1 variants would be pulled from Army stocks and refurbished for delivery by the fall. This change would also ensure deliveries to US allies of new M1A2s would not be disrupted. In September 2023, Ukraine began receiving these tanks, which were former U.S. Marine Corps tanks. The tanks supplied were also older (having entered service in 1986), "export" versions with classified US armor removed before the tanks were sent to Ukraine. This was to prevent or reduce any exploitation of technology found on any Abrams captured by Russian forces. In February 2024, an M1A1 was reported as lost in Ukraine. The blowout panels on the ammo bins had been activated, indicating that the ammunition had cooked off. This M1A1 was destroyed by a FPV Piranha 10 quadcopter. As of August 2024, Ukraine had visually confirmed losses of 14 (6 destroyed and 8 damaged and abandoned) of the 31 Abrams tanks, including one that was captured by Russia and displayed as a war trophy in Moscow in May 2024. One more Abrams was damaged. In April 2024, Pentagon officials reported that Ukraine's Abrams had been withdrawn from frontline service. The Russian use of hunter killer drones have made it "too difficult" to operate the tanks in the current battlefield with "muddy ground hindering manoeuvrability". A Ukrainian company has unveiled a new set of "anti-drone steel screens", which weighs "430 kg [approximately 948 pounds]". Designed to protect the tank, while not hindering its function, the screens also use Soviet era Kontakt-1 explosive reactive armor. The screens protect the turrets top, rear, sides and other vulnerable sections. It leaves opening for smoke grenade launchers, the commander's hatch and other parts of the tank. Some 7 sets of armor have been produced, according to the company, for the Ukrainian Abrams. In October 2024 Australia announced that 49 recently retired M1A1 tanks would be transferred to Ukraine as the Australian Army started receiving its new M1A2 models. Proposed production shutdown Serial production of the M1 Abrams for the U.S. Army ended in 1995, though production for exports continued until 2000. The U.S. Army planned to end operations at Joint Systems Manufacturing Center (formerly Lima Army Tank Plant) from 2013 to 2016 to save over $1 billion; it would be restarted in 2017 to upgrade existing tanks. General Dynamics Land Systems (GDLS), which operates the factory, opposed the move, arguing that suspension of operations would increase long-term costs and reduce flexibility. Specifically, GDLS estimated that closing the plant would cost $380 million and restarting production would cost $1.3 billion. By August 2013, Congress had allocated $181 million for buying parts and upgrading Abrams systems to mitigate industrial base risks and sustain development and production capability. Congress and General Dynamics were criticized for redirecting money to keep production lines open and accused of "forcing the Army to buy tanks it didn't need." General Dynamics asserted that a four-year shutdown would cost $1.1–1.6 billion to reopen the line, depending on the length of the shutdown, whether machinery would be kept operating, and whether the plant's components would be completely removed. They contended that the move was to upgrade Army National Guard units to expand a "pure fleet" and maintain production of identified "irreplaceable" subcomponents. A prolonged shutdown could cause their makers to lose their ability to produce them and foreign tank sales were not guaranteed to keep production lines open. There is still a risk of production gaps even with production extended through 2015. With funds awarded before recapitalization is needed, budgetary pressures may push planned new upgrades for the Abrams from 2017 to 2019. In December 2014, Congress again allocated $120 million, against the wishes of the Army, for Abrams upgrades. In late 2016, tank production and refurbishment had fallen to a rate of one per month with fewer than 100 workers on site. In 2017, the Presidency of Donald Trump ordered military production to increase, including Abrams production and employment. In 2018, it was reported that the Army had ordered 135 tanks rebuilt to new standards, with employment at over 500 workers and expected to rise to 1,000. The Marine Corps pursued a force restructuring plan named Force Design 2030. Under this program, all U.S. Marine tank battalions were deactivated and their M1A1 tanks transferred to the Army by the end of 2021. Future plans During the 1980s and 1990s, the Block III main battle tank from the Armored Systems Modernization (ASM) program was expected to succeed the M1 Abrams family in the 1990s. The design had an unmanned turret with a 140 mm main gun, as well as improved protection. The end of Cold War hostilities caused the end of the program. The tracked M8 Armored Gun System was conceived as a possible supplement for the Abrams in U.S. service for low-intensity conflict in the early 1990s. Prototypes were made but the program was canceled. The eight-wheeled M1128 mobile gun system was designed to supplement the Abrams in U.S. service for low-intensity conflicts. It has been introduced into service and serves with Stryker brigades. The Future Combat Systems XM1202 Mounted Combat System was to replace the Abrams in U.S. Army service and was in development when funding for the program was canceled in 2010. Engineering Change Proposal 1 is a two-part upgrade process. ECP1A adds space, weight, and power improvements and active protection against improvised explosive devices. Nine ECP1A prototypes have been produced as of October 2014. ECP1B, which would begin development in 2015, may include sensor upgrades and converging several tank round capabilities into a multipurpose round. As of 2021, the Army anticipated that the remaining M1A2 to beyond 2050. As of 2021 the Army is to begin divesting its M1A1 SA variants in fiscal year 2025. As of March 2023 the US Army had a stated goal of procuring 2,204 M1A2SEPv3 tanks with funds already having been committed to procure 2,093 of this variant. This will make the M1A2SEPv3 the standard issue tank for the US Army and US Army National Guard. As of 2021, the U.S. Army was evaluating a replacement for the M1 Abrams as part of the Next Generation Combat Vehicle (NGCV) program, notionally known as the Decisive Lethality Platform (DLP). In September 2023, the U.S. Army announced that it had canceled the planned M1A2 SEPv4 variant and would instead redirect resources into a new variant of the Abrams tank, named M1E3. Design Countermeasures Camouflage Some XM1 FSED pilot vehicles and XM1 LRIP tanks were painted with the Mobility Equipment Research and Design Command (MERDC) 4-color paint scheme. Factory-applied forest green paint gave way to "Europe 1", a three-color pattern, in 1983 at the same time as Chemical Agent Resistant Coating (CARC) was adopted. Europe 1 consisted of Green 383, Brown 383, and black colors. U.S Army Abrams deployed to the Iraq War were painted Carc Tan 686A. Due to the increasing significance of American operations in Europe, the U.S. Army transitioned most of its vehicles to CARC Green 383 starting around 2017. M1A1s came from the factory with the NATO three color camouflage Black/Med-Green/Dark-Brown CARC paint jobs. Today, M1A1s are given the NATO three color paint job during rebuilds. M1s and M1A1s deployed to Operation Desert Storm were hastily painted desert tan. Some, but not all, of these tanks were repainted to their "authorized" paint scheme. M1A2s built for Middle Eastern countries were painted in desert tan. Replacement parts (roadwheels, armor skirt panels, drive sprockets, etc.) are painted olive green, which can sometimes lead to vehicles with a patchwork of green and desert tan parts. Australian M1A1s are camouflaged in AUSCAM, a scheme that consists of black, olive drab, and brown. Concealment The turret is fitted with two six-barreled M250 smoke grenade launchers (USMC M1A1s used an eight-barreled version), with one on each side. When deployed, the grenades airburst, creating a thick smoke that blocks both visual and thermal imaging. The engine is also equipped with a vehicle engine exhaust smoke system (VEESS) that is triggered by the driver. When activated, fuel is sprayed into the hot turbine exhaust, creating thick smoke. This system was discontinued by the U.S. Army after it switched to JP-8 jet fuel in the 1990s due to the risk of fire. Armor In addition to conventional rolled homogeneous armor (RHA), the Abrams uses a secret British-developed Chobham composite armor. The M1 Abrams composite armor (referred to as "special armor" by the U.S. Army) is most substantial at the front of the hull, where it is at its thickest. The front of the hull is armored with composites. The Abrams turret features composite armor across both the front and the sides. The armor is much thicker on the Abrams than on previous tanks. This is not a reflection of any weakness of Chobham armor—pound-for-pound Chobham is better at stopping shaped charges and kinetic projectiles. Rather, unlike RHA, Chobham is optimized against shaped charge projectiles. Effective shaped charges, particularly anti-tank guided missiles, were a relatively new battlefield innovation. Lacking a breakthrough advance in novel armor material to negate shaped charges, previous tank designers had simply not found it practical to add the amount of RHA required to defeat shaped charges. While the exact composition of the Abrams' composite armor remains a state secret, a generalization about how it works can be gleaned from what has been publicly said about it. It consists of ceramic blocks set in resin between layers of conventional armor. The ceramic acts as a non-explosive reactive armor (NERA), disrupting shaped charges. The NERA plates shatter on impact with the projectile, disrupting the penetrating jets of shaped charges; or in the case of kinetic rounds eroding the projectile. For the M1 Abrams base model, military historian Steven Zaloga estimates the frontal armor at 350 mm vs APFSDS and 700 mm vs HEAT warhead in the book, M1 Abrams Main Battle Tank 1982–1992 (1993). In M1 Abrams vs T-72 Ural (2009), he uses Soviet estimates of vs APFSDS and vs HEAT for the base model Abrams. He also gives the Soviet estimates for the M1A1, vs APFSDS, and vs HEAT. Armor protection against kinetic energy rounds was improved by implementing a new special armor incorporating depleted uranium (DU). This was introduced into the M1A1 production starting October 1988. but at the expense of adding considerable weight to the tank, as depleted uranium is 1.7 times denser than lead. The DU is applied to the backing plate of the turret armor arrays. The first M1A1 tanks to receive this upgrade were tanks stationed in Germany. US-based tank battalions participating in Operation Desert Storm received an emergency program to upgrade their tanks with depleted uranium armor immediately before the onset of the campaign. M1A2 tanks uniformly incorporate depleted uranium armor, and all M1A1 tanks in active service have been upgraded to this standard as well. This variant was designated as the M1A1HA (HA for Heavy Armor). The M1A1 AIM, M1A2 SEP and all subsequent Abrams models feature depleted uranium. Each Abrams variant after the M1A1 have been equipped with depleted uranium armor of different generations. The M1A1HA uses first-generation armor, while the M1A2 and M1A1HC use second generation depleted uranium. The M1A2 SEP variants have been equipped with third-generation depleted uranium armor combined with a graphite coating. All U.S Army Abrams since 1998 received Depleted Uranium in the Hull, Turrets and Side of the turret, with the SEP Variants receiving third generation D.U. For the M1A1HA, Zaloga gives a frontal armor estimate of vs APFSDS and vs HEAT in M1 Abrams Main Battle Tank 1982–1992, nearly double the original protection of the Abrams. In M1 Abrams vs T-72 Ural, he uses different estimates of vs APFSDS and vs HEAT for the front hull and vs APFSDS and vs HEAT for the front of the turret. The protection of M1A2 SEP is a frontal turret armor estimate of vs APFSDS and vs HEAT, glacis estimate of vs APFSDS and vs HEAT, and lower front hull estimate of vs APFSDS and vs HEAT. The M1A2 SEPv3 increased the LOS thickness of the turret and hull front armor; total armor protection from this increase is not known. In 1998, a program was begun to incorporate improved hull, turret, and side armor into the M1A2. This was intended to offer better protection against rocket-propelled grenades that were more modern than the baseline RPG-7. These kits were installed on about 325 older M1A2 tanks in 2001–2009 and were also included in upgraded tanks. The Abrams may also be fitted with explosive reactive armor over the track skirts if needed (such as the Tank Urban Survival Kit) and slat armor over the rear of the tank and rear fuel cells to protect against ATGMs. The 105 mm M1 Abrams does not use spall liners, though three 105 mm rounds on the turret basket floor are covered with spall protection covers on the M1 tank variant. Damage control The tank has a halon firefighting system to automatically extinguish fires in the crew compartment. The engine compartment has a firefighting system that is engaged by pulling a T-handle located on the left side of the hull. The Halon gas can be dangerous to the crew. However, the toxicity of Halon 1301 gas at 7% concentration is much lower than the combustion products produced by fire in the crew compartment, and CO2 dump would be lethal to the crew. The crew compartment also contains small hand-held fire extinguishers. Fuel and ammunition are stored in armored compartments with blowout panels intended to protect the crew from the risk of the tank's own ammunition cooking off (exploding) if the tank is damaged. The main gun's ammunition is stored in the rear section of the turret, with blast doors that open under power by sliding sideways only to remove a round for firing, then automatically close. Doctrine mandates that the ammunition door must be closed before arming the main gun. NBC protection Starting with the M1A1 variant nuclear, biological, chemical protection was provided by a turret overpressure system. Previously the Abrams crew had been required to don NBC suits in case of an NBC attack. NBC masks are still retained as a backup, and crews often train while wearing them to remain proficient and combat-effective in such a scenario. Tank Urban Survival Kit The Tank Urban Survival Kit (TUSK) is a series of improvements to the M1 Abrams intended to improve fighting ability in urban environments. Historically, urban and other close battlefields have been poor places for tanks to fight. A tank's front armor is much stronger than that on the sides, top, or rear. In an urban environment, attacks can come from any direction, and attackers can get close enough to reliably hit weak points in the tank's armor or gain sufficient elevation to hit the top armor. Armor upgrades include reactive armor on the sides of the tank and slat armor on the rear to protect against rocket-propelled grenades and other shaped charge warheads. Abrams Reactive Armor Tile (ARAT) I consists of 32 XM19 reactive armor boxes added to the sides of the tank. ARAT II consists of rounded XM32 reactive armor tiles mounted over-top the XM19 tiles. A Transparent Armor Gun Shield and a thermal sight system are added to the loader's top-mounted M240B 7.62 mm machine gun, and a Kongsberg Gruppen Remote Weapon Turret carrying a 12.7 mm (.50 in) caliber machine gun (again similar to that used on the Stryker) is in place of the tank commander's original 12.7 mm (.50 in) caliber machine gun mount, wherein the commander had to expose himself to fire the weapon manually. An exterior telephone allows supporting infantry to communicate with the tank commander. In August 2006, General Dynamics Land Systems received a U.S. Army order for 505 Tank Urban Survivability Kits (TUSK) for Abrams main battle tanks supporting operations in Iraq, under a US$45 million contract. Deliveries were expected to be completed by April 2009. Under a separate order, the U.S. Army awarded General Dynamics Armament and Technical Products (GDATP) US$30 million to produce reactive armor kits to equip M1A2s. Tiles will be produced at the company's reactive armor facility in Stone County Operations, McHenry, Mississippi. In December 2006, the U.S. Army added Counter Improvised Explosive Device enhancements to the M1A1 and M1A2 TUSK, awarding GDLS $11.3 million contract, part of the $59 million package mentioned above. In December, GDLS also received an order, amounting to around 40% of a US$48 million order, for loader's thermal weapon sights being part of the TUSK system improvements for the M1A1 and M1A2 Abrams Tanks. Active protection system In addition to the armor, some USMC Abrams tanks were equipped with a soft-kill active protection system, the AN/VLQ-6 Missile Countermeasure Device (MCD) that can impede the function of guidance systems of some semi-active control line-of-sight (SACLOS) wire- and radio guided anti-tank missiles (such as the Russian 9K114 Shturm) and infrared homing missiles. These were not ready in time for the Gulf War. The MCD works by emitting a massive, condensed infrared signal to confuse the infrared homing seeker of an anti-tank guided missile (ATGM). However, the drawback to the system is that the ATGM is not destroyed, it is merely directed away from its intended target, leaving the missile to detonate elsewhere. During the Iraq War the U.S. Marine Corps equipped its M1A1s with AN/VLQ-8A electro-optical jammers. In 2016, the U.S. Army and Marine Corps began testing the Israeli Trophy active protection system to protect their Abrams tanks from modern RPG and ATGM threats by either jamming (with ATGMs) or firing small rounds to deflect incoming projectiles. The Army planned to field a brigade of over 80 tanks equipped with Trophy to Europe in 2020. It is planned for up to 261 Abrams to be upgraded with the system, enough for four brigades. In June 2018, the Army awarded Leonardo DRS, U.S. partner to Trophy's designer Rafael, a $193 million contract to deliver the system in support of M1 Abrams "immediate operational requirements". U.S. Army M1A2 SEPv2 Abrams tanks deployed to Germany in July 2020 fitted with Trophy systems. Deliveries to equip four tank brigades were completed in January 2021. Armament Primary M68A1 rifled gun The main armament of the original model M1 and IPM1 was the M68A1 105 mm rifled tank gun firing a variety of APFSDS, HEAT, high explosive, white phosphorus rounds and an anti-personnel (multiple flechette) round. This gun used a license-made tube of the British Royal Ordnance L7 gun together with the vertical sliding breech block and other parts of the U.S. T254E2 prototype gun. However, a longer ranged weapon was always envisaged, with lethality beyond to combat newer armor technologies. To attain that lethality, the projectile diameter needed to be increased. The tank was able to carry 55 105 mm rounds, with 44 stored in the turret blowout compartment and the rest in hull stowage. Being non-combustible, the empty cartridge cases of the M1 variant accumulated on the turret floor after firing. After allowing some time to cool, they were ejected out of the hatch by the loader. M256 smoothbore gun The main armament of the M1A1 and M1A2 is the M256 120 mm smoothbore gun, designed by Rheinmetall AG of Germany, manufactured under license in the U.S. by Watervliet Arsenal, New York. The M256 is an improved variant of the Rheinmetall 120 mm L/44 gun carried on the German Leopard 2 on all variants up to the Leopard 2A5, the difference being in thickness and chamber pressure. Leopard 2A6 replaced the L/44 barrel with a longer L/55. Due to the increased caliber, only 40 or 42 rounds are able to be stored depending on if the tank is an A1 or A2 model. Elevation: −9 to +20 degree The M256 fires ammunition with combustible cartridge cases made out of nitrocellulose. The cartridges were safer against premature ignition and flarebacks than earlier combustible cartridge rounds, but not entirely accident-proof. The M256 fires a variety of rounds. The primary APFSDS round of the Abrams is the depleted uranium M829 round, of which four variants have been designed. M829A1, known as the "Silver Bullet", saw widespread service in the Gulf War, where it proved itself against Iraqi armor such as the T-72. The M829A2 APFSDS round was developed specifically as an immediate solution to address the improved protection of a Russian T-72, T-80U or T-90 main battle tank equipped with Kontakt-5 explosive reactive armor (ERA). Later, the M829A3 round was introduced in 2002 to improve its effectiveness against next-generation ERA equipped tanks. Development of the M829 series is continuing with the M829A4 currently entering production, featuring advanced technology such as data-link capability. The Abrams also fires HEAT warhead shaped charge rounds such as the M830, the latest version of which (M830A1) incorporates a sophisticated multi-mode electronic sensing fuse and more fragmentation that allows it to be used effectively against armored vehicles, personnel, and low-flying aircraft. The Abrams uses a manual loader, who also provides additional support for maintenance, observation post/listening post (OP/LP) operations, and other tasks. The new M1028 120 mm anti-personnel canister cartridge was brought into service early for use in the aftermath of the 2003 invasion of Iraq. It contains 1,098 tungsten balls that spread from the muzzle to produce a shotgun effect lethal out to . The tungsten balls can be used to clear enemy dismounts, break up hasty ambush sites in urban areas, clear defiles, stop infantry attacks and counter-attacks and support friendly infantry assaults by providing covering fire. The canister round is also a highly effective breaching round and can level cinder block walls and knock man-sized holes in reinforced concrete walls for infantry raids at distances up to . Also in use is the M908 obstacle-reduction round. It is designed to destroy obstacles and barriers. The round is a modified M830A1 with the front fuse replaced by a steel nose to penetrate into the obstacle before detonation. The U.S. Army Research Laboratory (ARL) conducted a thermal analysis of the M256 from 2002 to 2003 to evaluate the potential of using a hybrid barrel system that would allow for multiple weapon systems such as the XM1111 Mid-Range munition, airburst rounds, or XM1147. The test concluded that mesh density (number of elements per unit area) impacts accuracy of the M256 and specific densities would be needed for each weapon system. In 2013, the Army was developing a new round to replace the M830/M830A1, M1028, and M908. Called the M1147 Advanced Multi-Purpose (AMP) round, it will have point detonation, delay, and airburst modes through an ammunition data-link and a multi-mode, programmable fuse in a single munition. Having one round that does the job of four would simplify logistics and be able to be used on a variety of targets. The AMP is to be effective against bunkers, infantry, light armor, and obstacles out to , and will be able to breach reinforced concrete walls and defeat ATGM teams from . Orbital ATK was awarded a contract to begin the first phase of development for the AMP XM1147 High-Explosive Multi-Purpose with Tracer cartridge in October 2015. As of 2024 the round is undergoing the final testing stages, with the full-rate production decision scheduled for the end of the year. In addition to these, the XM1111 (Mid-Range-Munition Chemical Energy) was also in development. The XM1111 was a guided munition using a dual-mode seeker that combined imaging-infrared and semi-active laser guidance. The MRM-CE was selected over the competing MRM-KE, which used a rocket-assisted kinetic energy penetrator. The CE variant was chosen due to its better effects against secondary targets, providing a more versatile weapon. The Army hoped to achieve IOC with the XM1111 by 2013. However, the Mid-Range Munition was canceled in 2010 along with Future Combat Systems. Secondary The Abrams tank has three machine guns, with an optional fourth: A .50 cal. (12.7 mm) M2HB machine gun in front of the commander's hatch. On the M1 and M1A1, this gun is mounted on the Commander's Weapons Station. This allows the weapon to be aimed and fired from within the tank. Normal combat loadout for the M1A1 is a single 100-round box of ammo at the weapon, and another 900 rounds carried. The later M1A2 variant had a "flex" mount that required the tank commander to expose his or her upper torso in order to fire the weapon. In urban environments in Iraq this was found to be unsafe. With the Common Remote Operated Weapons System (CROWS) add-on kit, an M2A1 .50 Caliber Machine gun, M240, or M249 SAW can be mounted on a CROWS remote weapons platform (similar to the Protector M151 remote weapon station used on the Stryker family of vehicles). Current variants of the Tank Urban Survival Kit (TUSK) on the M1A2 have forgone this, instead adding transparent gun shields to the commander's weapon station. The upgrade variant called the M1A1 Abrams Integrated Management (AIM) equips the .50 caliber gun with a thermal sight for accurate night and other low-visibility shooting. A 7.62 mm M240 machine gun in front of the loader's hatch on a skate mount (seen at right). Some of these were fitted with gun shields during the Iraq War, as well as night-vision scopes for low-visibility engagements and firing. This gun can be moved to the TC's position if the M2 .50 cal is damaged. A second 7.62 mm M240 machine gun in a coaxial mount (i.e., it points at the same targets as the main gun) to the right of the main gun. The coaxial MG is aimed and fired with the same computerized firing control system used for the main gun. On earlier M1 and M1A1s 3000 rounds are carried, all linked together and ready to fire. This was reduced slightly in later models to make room for new system electronics. A typical 7.62 mm combat loadout is between 10,000 and 14,000 rounds carried on each tank. (Optional) A second coaxial .50 cal. (12.7 mm) M2HB machine gun can be mounted directly above the main gun in a remote weapons platform as part of the CSAMM (Counter Sniper Anti Material Mount) package. Aiming The Abrams is equipped with a ballistic fire-control computer that uses user and system-supplied data from a variety of sources to compute, display, and incorporate the three components of a ballistic solution—lead angle, ammunition type, tube wear, propellant temperature, wind speed, air temperature, the relative motions of the target and the Abrams, and range to the target—to accurately fire the main gun. These three components are determined using a laser rangefinder, crosswind sensor, a pendulum static cant sensor, data concerning performance and flight characteristics of each specific type of round, tank-specific boresight alignment data, ammunition temperature, air temperature, barometric pressure, a muzzle reference system (MRS) that determines and compensates for barrel drop at the muzzle due to gravitational pull and barrel heating due to firing or sunlight, and target speed determined by tracking rate tachometers in the Gunner's or Commander's Controls Handles. All of these factors are computed into a ballistic solution and updated 30 times per second. The updated solution is displayed in the Gunner's or Tank Commander's field of view in the form of a reticle in both day and thermal modes. The ballistic computer manipulates the turret and a complex arrangement of mirrors so that all one has to do is keep the reticle on the target and fire to achieve a hit. Proper lead and gun tube elevation are applied to the turret by the computer, greatly simplifying the job of the gunner. The fire control system on the M1 and M1A1 variants is the Computing Devices Canada ballistic computer system. On the M1A2 the Fire Control Electronics Unit is manufactured by GDLS. The laser designator is a Hughes model. This fire control system uses this data to compute a firing solution for the gunner. The ballistic solution generated ensures a hit percentage greater than 95 percent at nominal ranges. Either the commander or gunner can fire the main gun. Additionally, the Commander's Independent Thermal Viewer (CITV) on the M1A2 can be used to locate targets and pass them on for the gunner to engage while the commander scans for new targets. If the primary sight system malfunctions or is damaged, the main and coaxial weapons can be manually aimed using a telescopic scope boresighted to the main gun known as the Gunner's Auxiliary Sight (GAS). The GAS has two interchangeable reticles; one for HEAT and multi-purpose anti-tank (MPAT) ammunition and one for APFSDS and Smart Target-Activated Fire and Forget (STAFF) ammunition. Turret traverse and main gun elevation can be performed with manual handles and cranks if the fire control or hydraulic systems fail. The commander's M2HB .50 caliber machine gun on the M1 and M1A1 is aimed by a 3× magnification sight incorporated into the Commander's Weapon Station (CWS), while the M1A2 uses the machine gun's own iron sights, or a remote aiming system such as the Common Remotely Operated Weapon Station (CROWS) system when used as part of the Tank Urban Survival Kit. The loader's M240 machine gun is aimed either with the built-in iron sights or with a thermal scope mounted on the machine gun. In late 2017, the 400 USMC M1A1 Abrams were to be upgraded with better and longer-range sights on the Abrams Integrated Display and Targeting System (AIDATS) replacing the black-and-white camera view with a color sight and day/night thermal sight, simplified handling with a single set of controls, and a slew to cue button that repositions the turret with one command. Preliminary testing showed the upgrades reduced target engagement time from six seconds to three by allowing the commander and gunner to work more closely and collaborate better on target acquisition. Mobility Tactical The M1 Abrams's powertrain consists of an AGT1500 multifuel gas turbine (originally made by Lycoming, now Honeywell) capable of at 30,000 rpm and at 10,000 rpm and a six-speed (four forward, two reverse) Allison X-1100-3B Hydro-Kinetic automatic transmission. This gives it a governed top speed of on paved roads, and cross-country. With the engine governor removed, speeds of around are possible on an improved surface. However, damage to the drivetrain (especially to the tracks) and an increased risk of injuries to the crew can occur at speeds above . The tank was built around this engine and it is multifuel-capable, including diesel, gasoline, marine diesel and jet fuel (such as JP-4 or JP-8). In the AGT1500, jet fuel has poorer fuel economy and operating range compared to diesel. By 1989, the Army was transitioning solely to JP-8 for the M1 Abrams, part of a plan to reduce the service's logistics burden by using a single fuel for aviation and ground vehicles. However, as of 2023, the U.S. Army frequently refuels the Abrams with diesel, which is also used by the Bradley Fighting Vehicle. The Australian M1A1 AIM SA burns diesel fuel, since the use of JP-8 is less common in the Australian Army. The gas turbine propulsion system has proven quite reliable in practice and combat, but its high fuel consumption is a serious logistic problem. It burns between 1.5 to 3 gallons per mile. The turbine is very quiet when compared to diesel engines of similar power output and produces a high-pitched whine, reducing the audible distance of the sound, thus earning the Abrams the nickname "whispering death" during its worldwide debut at the 1982 Reforger exercise. By the time production of the AGT1500 ended in 1994, the U.S. had purchased 12,000 such engines. In 2006 the Army awarded Honeywell a contract to overhaul 1000 engines, with options for up to 3000 more. The Army received proposals, including two diesel options, to provide the common engine for the XM2001 Crusader and Abrams. In 2000, the Army selected the gas turbine engine LV100-5 from Honeywell and subcontractor General Electric. The new LV100-5 engine was smaller (43% fewer parts) with rapid acceleration, quieter running, and no visible exhaust. It also featured a 33% reduction in fuel consumption (50% less when idle) and near drop-in replacement. The Common Engine Program was shelved when the Crusader program was canceled. Phase 2 of Army's PROSE (Partnership for Reduced O&S Costs, Engine) program, however, called for further development of the LV100-5 and replacement of the current AGT1500 engine. From 1991 to 1994, the Army fitted 1,500 Abrams turrets with external auxiliary power units (APU). APUs allow some the Abrams to run some functions without running on the engine. Some Abrams tanks that saw service during the Gulf War were fitted with such a device. Although the Army favored an under-armor APU, Congress instead funded a short-term modification to 336 M1A2 Abrams. These were installed in 1997. An under-armor APU located in the hull was chosen for the M1A2 SEP variant. When this proved unreliable, it was replaced with a battery-based Alternate APU starting in 2005. Although the M1 tank is not designed to carry riders easily, provisions exist for the Abrams to transport troops in tank desant with the turret stabilization device switched off. A battle-equipped infantry squad may ride on the rear of the tank, behind the turret. The soldiers can use ropes and equipment straps to provide handholds and snap links to secure themselves. The Abrams T156 is a permanently bonded rubber track pad, a distinctive feature not found on any other tank. Unlike other tanks with replaceable track pads, on the Abrams, a worn track pad is remedied by replacing the entire track shoe. The Abrams non-removable track pads save weight but are less desirable in snow as the pads cannot be replaced with grousers. As of 2007, M1 Abrams track wear constitutes the second-largest consumable expense in the U.S. Army, surpassed only by Meals, Ready to Eat consumption. In 1988 the Army awarded FMC Corporation a contract for T158 tracks rated for , or about double the life of the previous shoe. These feature replaceable pads and are about 3000 pounds heavier. The driver is equipped with a thermal viewer. On at least some models this is the Hughes AN/VAS-3. Strategic Strategic mobility is the ability of the tanks of an armed force to arrive in a timely, cost effective, and synchronized fashion. The Abrams can be carried by a C-5 Galaxy or a C-17 Globemaster III. The limited capacity (two combat-ready tanks in a C-5, one combat-ready tank in a C-17) caused serious logistical problems when deploying the tanks for the first Gulf War, though there was enough time for 1,848 tanks to be transported by ship. The Marines transported their Marine Air-Ground Task Force Abrams tanks by combat ship. A Wasp-class Landing Helicopter Dock (LHD) typically carried a platoon of four to five tanks attached to the deployed Marine Expeditionary Unit, which were then amphibiously transported to shore by Landing Craft Air Cushion (LCAC) at one combat-ready tank per landing craft. The Abrams is also transportable by truck, namely the Oshkosh M1070 and M1000 Heavy Equipment Transporter System (HETS) for the US Military. The HETS can operate on highways, secondary roads, and cross-country. It accommodates the four tank crew members. The Australian Army uses customized MAN trucks to transport its Abrams. The first instance of the Abrams being airlifted directly into a battlefield occurred in October 1993. Following the Battle of Mogadishu, 18 M1 tanks were airlifted by C-5 aircraft to Somalia from Hunter Army Airfield, Georgia. Issues Air filter clog In a NSIA report on the Abrams in the Gulf War, crews reported issues related to the turbine engine, other than the fuel consumption concerns, they noted the Abrams suffered from sand clogging the filters which were known to cause reduced fuel economy, or in the worst case, engine damage. Doctrine, crew responsibilities and platoon operations Before the M1 Abrams program, the U.S. Army had designed tanks to conform to doctrine. This approach changed with the XM1, where the Army wrote its doctrine after developing the tank. The U.S. Army's Abrams tank doctrine was influenced by German, British, American, and Soviet ideas. The German concept of Auftragstaktik (English: Mission-type tactics), a military doctrine emphasizing decentralized decision-making, and Schwerpunkt (English: Main point), the massing of resources around a focal point, were influential. German-type breakthrough tactics favored by general George S. Patton, commander of the Seventh Army in the Mediterranean Theater of World War II were advocated by Creighton Abrams and his devotees U.S. Army Training and Doctrine Command (TRADOC) commanding generals William E. DePuy and his successor Donn A. Starry. The Army's new fighting doctrine was drafted by TRADOC commanding general DePuy, and was heavily influenced by lessons from the 1973 Arab–Israeli Yom Kippur War. Field Manual 100-5 Operations, published in 1976, became "one of the most controversial documents the Army had ever published", according to Orr Kelly. The document recognized that U.S. forces would quickly become outnumbered in the case of a surprise Soviet invasion. It called for U.S. forces to maneuver quickly to where they were needed to mount an "active defense" oriented towards blunting the spearhead of the attacking force. Critics of this document noted that Soviet attacks would come in waves that would overwhelm U.S. defenses. The revision to the manual, which faced criticism rivaling that of the first edition, was published in 1982. The manual's emphasis was influenced by Depuy's successor, Starry. It called for using the "entire depth of the battlefield to strike the enemy and to prevent him from concentrating his firepower or manuevering his forces to a point of his choice." This alarmed NATO allies, who considered U.S. counterattacks across enemy borders to be needlessly provocative. The third revision of the manual published in 1986 left open the possibility of attacking across enemy borders at the discretion of politicians. When the Abrams entered U.S. Army service in 1980, its arrival marked an organizational change. The tank battalion went from three companies of three platoons to four companies of three platoons. The standard tank platoon fell from five tanks — a number consistent since the first days of the Tank Corps in World War I — to four. The change reflected both the improved capability of the new tank but also its cost. The reduction in platoon size necessitated changes in tactics oriented upon platoon and section actions in which the platoon leader had both to fight his tank and manage the unit. United States Platoon organization within the U.S. Army and U.S. Marine Corps as of 2019 is as follows: A tank platoon includes four Abrams MBTs organized into two sections, with two tanks in each section. "A" section consists of the platoon leader (platoon commander in USMC parlance) who is the commander of the vehicle designated as Tank 1, and the platoon leader's wingman, who is the commander of Tank 2. "B" section consists of the platoon sergeant, who is the tank commander of Tank 4, and Tank 3 is the platoon sergeant's wingman. The wingman concept requires that individual tanks orient off the tank to its left or right side. In the tank platoon, Tank 2 orients off the platoon leader's tank, while Tank 3 orients off the platoon sergeant's tank. The platoon sergeant orients off the platoon leader's tank. The tank platoon is organic to Armor companies of a combined arms battalion. The platoon may be attached to a number of organizations, commonly a mechanized infantry company, to create company teams. It may also be placed under the control of an Infantry organization. The exact amount of control the gaining unit would have is determined by the command relationship established by its higher HQ. The Armor company is organized, equipped, and trained to fight pure or as a task organized company team. The Armor company includes an HQ and three tank platoons. The company headquarters is equipped with two MBTs, armored personnel carriers, and wheeled vehicles for mission command/command and control and sustainment. Maintenance and Operation A series of TM 9 technical manuals cover various aspects of the tanks maintenance and operation. The exact number and titles of TM 9 manuals for the M1 Abrams may vary depending on the specific variant (M1, M1A1, etc.) and the date of publication. M1: Initial production model with a 105mm gun. M1A1: Upgraded with a 120mm smoothbore gun, improved armor, and a bustle rack. M1A2: Features a commander's independent thermal viewer (CITV), improved fire control systems, and an enhanced digital architecture. M1A2 SEP (System Enhancement Package): A series of upgrades to the M1A2, including SEP v2 and SEP v3, with further improvements. Hull Engine and Powertrain Engine: type, its components, and maintenance procedures. Transmission: system type, including gearboxes and differentials. Air Intake and Filtration System: components and their function. Mobility Systems Suspension: system type, including road wheels and bogies. Tracks: type and replacement procedure. Steering and Braking System: type and their operation. Ukraine A Ukrainian tank battalion consists of 31 tanks: three companies of three platoons, with each tank platoon operating three tanks. Each company and battalion is led by an additional command tank. Variants and upgrades General Motors XM1 validation phase prototype. Chrysler XM1 validation phase prototype. XM1-FSED: Chrysler preproduction test model. Eleven Full-Scale Engineering Development test bed vehicles were produced in 1977–78. These vehicles were also called Pilot Vehicles and numbered PV-1 through PV-11. M1: First production variant. Production began in 1979 and continued to 1985 (3,273 built for the United States). The first 110 tanks were low rate initial production (LRIP) models, still called XM1s, because they were built before the tank was type-classified as the M1. The M1 variant was retired from active U.S. Army service in 1996. IPM1 (Improved Performance): Produced from 1984 to 1986 concurrent to the M1A1. It contained upgrades and reconfigurations like a rear storage bustle rack, improved armor, suspension, transmission, and final drives. (894 built for United States). : Production started in 1985 and continued to 1992, pressurized NBC system, rear bustle rack for improved stowage of supplies and crew belongings, redesigned blowoff panels and M256 120 mm smoothbore cannon (4,976 built for the U.S. Army, 221 for USMC, 59 M1A1 AIM SA sold to Australia). M1A1HA (Heavy Armor): Added first generation depleted uranium armor components. Some tanks were later upgraded with second generation depleted uranium armor components, and are unofficially designated M1A1HA+. M1A1HC (Heavy Common): Added new second generation depleted uranium armor components, digital engine control and other small upgrades common between Army and Marine Corps tanks. M1A1D (Digital): A digital upgrade for the M1A1HC, to keep up with M1A2 SEP, manufactured in quantity for only 2 battalions. M1A1 AIM v.1 (Abrams Integrated Management): A program whereby older units are reconditioned to initial factory standards, and the tank is improved by adding Forward-Looking InfraRed (FLIR) and Far Target Locate sensors, a tank-infantry phone, communications gear, including FBCB2 and Blue Force Tracking to aid in crew situational awareness, and a thermal sight for the .50 caliber machine gun. M1A1 AIM v.2/M1A1SA (Situational Awareness): Upgrades similar to AIM v.1 tanks and new third generation depleted uranium armor components. Configuration for the Royal Moroccan Army, which is almost identical to the Australian variant, except exportable turret armor is installed by General Dynamics Land System to replace the DU armor. M1A1 SA-UKR: Official U.S. designation for M1A1SA variants given to Ukraine via Foreign Military Sales program. M1A1 FEP (Firepower Enhancement Package): Similar upgrade to AIM v.2 for USMC tanks. M1A1KVT (Krasnovian Variant Tank): M1A1s that have been visually modified to resemble Soviet-made tanks for use at the National Training Center, fitted with MILES gear and a Hoffman device. M1A1M: An export variant ordered by the Iraqi Army. M1A1 (AIDATS upgrade): Upgrade-only variant to all USMC General Dynamics M1A1 Abrams tanks to improve the tank commander's situational awareness with an upgraded thermal sight, color day camera, and a stationary color display. (Baseline): Production began in 1992 and initial operating capability achieved in 1993. (77 built for the U.S. and more than 600 M1s upgraded to M1A2, 315 for Saudi Arabia, 1,005 for Egypt, 218 for Kuwait). The M1A2 offers the tank commander an independent thermal sight and ability to, in rapid sequence, shoot at two targets without the need to acquire each one sequentially, also second-generation depleted uranium armor components. M1A2 SEP (System Enhancement Package): Is fitted with new, second-generation gunner's thermal sight. Has upgraded third-generation depleted uranium armor components with graphite coating (240 new built, 300 M1A2s upgraded to M1A2 SEP for the United States, also unknown numbers of upgraded basic M1s and IPM1s, also 400 oldest M1A1s upgraded to M1A2 SEP). M1A2S (Saudi Arabian Package): Saudi Arabian variant upgrade of the M1A2 based on M1A2 SEP, with some features, such as depleted uranium armor, believed to be missing and replaced by special armor. (442 M1A2s upgraded to M1A2S). M1A2 SEPv2: Added Common Remotely Operated Weapon Station low-profile as standard, color displays, improved interfaces, a new operating system, improved crew-compartment cooling, and new second generation thermal optics. M1A2 SEPv3 (formerly M1A2C): Has increased power generation and distribution, better communications and networking, new Vehicle Health Management System (VHMS) and Line Replaceable Modules (LRMs) for improved maintenance, an Ammunition DataLink (ADL) to use airburst rounds, improved counter-IED armor package, Next Generation Armor Package (NGAP), and an Auxiliary Power Unit (APU) under armor to run electronics while stationary instead of the engine, visually distinguishing the version by a small exhaust at the left rear. Lethality enhancements include the M829A4 kinetic-energy anti-tank round enhance the tank's lethality against modern threats, such as advanced explosive reactive armor (AERA) and Active Protection Systems (APS). As well as, the M1147 AMP round that combines multiple functionalities, including point detonate, delay, and airburst modes, replacing four older round types and providing capabilities for obstacle reduction, bunker defeat, and precision airburst against anti-tank missile teams. The Ammunition DataLink (ADL) enables the round to communicate with the fire-control system, allowing the crew to program the desired mode in real-time for maximum effectiveness. The SEPv3 also has Improved Forward-Looking Infrared (IFLIR) technology, which significantly improves target acquisition, identification, and engagement under all conditions, including obscurants such as fog or smoke. The IFLIR integrates long-wave and mid-wave infrared sensors into both the gunner’s primary sight and the commander’s independent thermal viewer, offering enhanced detection capabilities at greater ranges. It provides four fields of view (FOV) displayed on high-definition screens, enabling faster and more accurate engagement of targets compared to the older second-generation FLIR systems. The Low-Profile CROWS (LP CROWS) significantly reduces the weapon station's profile, enhancing the tank commander’s situational awareness with upgraded day cameras featuring picture-in-picture technology, a 340% larger field of view in it's wide FOV mode, and improved targeting capabilities under both open- and closed-hatch conditions. More passive ballistic protection was added to the turret faces, along with new Explosive Reactive Armor mountings (Abrams Reactive Armor Tile (ARAT)) and Trophy Active Protection systems added to the turret sides. Prototypes for the SEPv3 began testing in 2015. The US Army is able to produce a maximum of 35 M1A2SEPv3 a month at the Lima plant in Ohio with a standard rate of 12 per month and 1 shift at 40 hours per week. The Army is producing the tank at a rate of 109 a year or roughly 9 a month. M1A2T: Special configuration variant of the M1A2 SEPv3 reportedly being offered for sale to Taiwan as of March 2019 and approved by U.S. Department of State as of July 2019. Per DSCA statement, it is roughly equivalent to M1A2 SEPv3, except depleted uranium armor is replaced by FMS export armor. There is no mention of the Trophy APS system. The new-built tanks will be produced at Anniston Army Depot, Anniston, Alabama, and the Joint Systems Manufacturing Center, Lima, Ohio. M1A2R: Variant for the Romanian Army, is under development. According to the chief of the General Directorate for Armaments, the Romanian Abrams variant will be a configuration of the M1A2 SEPv3. M1A2K: Variant for the Kuwaiti Army, slated to replace Kuwait's current M1A2 fleet. M1A2 SEPv4 (formerly M1A2D, canceled): The 2-12 Cavalry Regiment received the first M1A2 SEPv4 tank . It was previously under development . The Commander's Primary Sight, also known as the Commander's Independent Thermal Viewer, and Gunner's Primary Sight will be upgraded with third Gen FLIR, an improved laser rangefinder and color cameras. Additional improvements will include advanced meteorological sensors, laser warning/detection receivers, directional smoke grenade launchers and integration of the new XM1147 (AMP) 120 mm tank round. The AN/VVR-4 laser warning receiver and ROSY rapid obscurant system have been trialed by the US Army for adoption on the Abrams tank and Bradley fighting vehicle. The M1A2 SEPv4 variant was officially canceled by the U.S. Army on 6 September 2023. The Army plans to develop a new variant of the Abrams, designated M1E3, to enter service in the 2030s. M1E3: On 6 September 2023, the U.S. Army announced that it had canceled the planned M1A2 SEPv4 variant and would instead redirect resources into a new variant of the Abrams tank, named M1E3. The new variant is to include modular open-systems architecture and is designed to be lighter and more survivable on the battlefield. This variant is expected to be designated "M1A3" upon entering operational service. The Army Science Board report "An Independent Assessment of the 2040 Battlefield and its Implications for the 5th Generation Combat Vehicle (5GCV)" which reportedly influenced senior Army leadership to establish the program, recommending a $2.9 billion, eight/seven-year program to develop a "fifth generation combat vehicle," with proposed capabilities including: a hybrid electric drive an autoloader and new main gun advanced munitions, such as maneuvering hypersonic and gun-launched anti-tank guided missiles integrated armor protection improved command, control, and networking capabilities artificial intelligence (AI) applications; ability to pair with robotic vehicles; and masking capabilities to reduce the vehicle's thermal and electromagnetic signatures. M1 Abrams Block III Tank Test Bed (M1 TTB) was a prototype built in 1983 as part of TACOM's Abrams Block III program (whose purview was to eventually create the M1A3), featuring an unmanned turret with a 44-caliber 120 mm M256 smoothbore gun, three crew members sitting side by side inside an armored capsule at the front of the hull and a suite of cameras and thermal viewers to preserve the crew's situational awareness. The main armament was linked to a Meggitt mechanical autoloader and a 44-round vertical ammunition carousel inside the turret basket; this system could provide a sustained rate of fire of 10 rounds per minute and successfully completed over 40,000 round loading/unloading cycles without malfunctioning during testing. Due to the absence of a full-fledged manned turret and the lack of internal armor packages, the vehicle only weighed 45 tons. The program was cancelled after the end of the Cold War and its only prototype is now on display at the U.S. Army Armor & Cavalry Collection at Fort Moore, Georgia. M1 Thumper (also known as ATAC System Demonstrator) was a single M1A1 fitted with a heavily modified unarmored M1A1 turret to trial the experimental XM291 ATAC (sometimes referred to as LW120) smoothbore gun, a more powerful replacement for the M256 capable of firing either single-piece 120 mm or two-piece 140 mm ammunition with only a barrel change. The 140 mm rounds were too large (boasting twice the chamber volume of a M829 APFSDS and twice the muzzle energy) and heavy to be moved around by a human loader, mandating the installation of a XM91 mechanical cassette autoloader. The Thumper underwent testing in 1988 and in the 1990s at Aberdeen Proving Ground, where it demonstrated accuracy equal to an M1A1's but with significantly higher armor penetration capability. Component Advanced Technology Test Bed (CATTB) was a pair of highly modified prototypes conceived under the auspices of the Advanced Tank Armament Systems (ATAS) program to test several promising technologies. Two vehicles were built in 1993 and 1994: the first one, dubbed Phase I, paired a spare M1 hull with a modified turret (an evolution of the M1 Thumper's) fitted with a 48-caliber 140 mm Watervliet Arsenal XM291 smoothbore cannon and a Benét Labs XM91 mechanical autoloader inside the bustle; the second vehicle, codenamed Phase II, used the same turret as basis, but mated it to a brand new M1 hull altered to contain a more compact Cummins XAP-1000 AIPS diesel engine and two vertically-stacked, horizontal carousels (for non-ready ammunition) between the turret basket and the powerpack compartment. Phase II also trialed new single-shoe XT166 tracks (Phase I retained the original, two-shoed T156 of the M1), an in-arm hydropneumatic suspension and the Multi-Sensor Target Acquisition System (MTAS) with its low-power, millimetric wave radar. Both vehicles were tested extensively at Aberdeen Proving Ground. The aforementioned Thumper is often erroneously described as a de-tuned CATTB, although it predates the latter by five years. AbramsX is a technology demonstrator of the M1 Abrams series by General Dynamics Land Systems. The AbramsX features a lightweight Watervliet Arsenal XM360 smoothbore gun with pepperpot muzzle brake, an autoloader which reduces the crew to three, an unmanned turret, a hybrid diesel-electric Cummins ACE power pack that gives 50% more fuel efficiency, a 30 mm XM914 chain gun in a remote weapon station, Trophy active protection systems with three launchers, an augmented reality suite that would increase the crews' situational awareness thanks to cameras and sensors mounted around the tank's exterior, a silent mode when running on electric power, the ability to be updated more easily than existing tanks and use loitering munitions such as the AeroVironment Switchblade as well as surveillance drones, and reduced weight for improved mobility. In October 2022, GDLS released a video showing the Technology Demonstrator and various technology tests. Specialized Surrogate Research Vehicle: The surrogate research vehicle (SRV) project was conducted from 1980 to 1987 to evaluate the effectiveness of different crew arrangements using a turretless Abrams test bed. These modifications included adding two crew stations to the front hull. Two crew positions were retained in a rotating basket where the turret had been. Armored Recovery Vehicle: Initially known as the RV90, this was a prototype designed by General Dynamics. It was produced under contract with TACOM in 1988 despite an earlier preliminary decision to procure the M88A1E1. The prototype was evaluated against the M88A1E1 later that year. The Abrams RV was based on the Abrams chassis, but housed a crew of three in a unique armored superstructure. The Army selected the M88A1E1 regardless, which went into production as the M88A2 Hercules. Air Ground Defense System (AGDS): Proposed air defense variant of the Abrams equipped with dual 35 mm Bushmaster III autocannons, 12 ADATS missiles and advanced electro-optical and radar targeting systems derived from the ADATS. It was supposed to be capable of both air defense and anti-tank purposes with the ADATS MIM-146 missiles which was a dual purpose ATGM/SAM. The proposal never saw consideration and was never developed further. M1 Panther II: A mine-clearing vehicle with turret removed, mine rollers on the front, and magnetized dog bone. This could be operated remotely or with a crew of two. Six were built and two were deployed in 2007 by the USMC in Iraq. M104 Wolverine Heavy Assault Bridge: GDLS produced these under contract for the U.S. Army with testing beginning in 1996. The bridge was produced by Krauss-Maffei Wegmann. The chassis is an Abrams converted to M1A2 standard. Forty-three units were produced when production wrapped up in 2003. Battle Command Vehicle: The vehicle was visually modified with a mock M256 gun to appear like an ordinary Abrams MBT, but featured communications equipment and workstations for battle commanders. United Defense LP constructed a prototype which the Army tested at Fort Hood in 1997. Visually modified: The National Training Center possesses 28 visually modified M1A1s resembling T-80s. M1 Grizzly (Breacher): In 1997 TACOM awarded United Defense a $129 million contract to construct this vehicle based on the Abrams chassis. This was capable of clearing minefields and demolishing obstacles with its dozer blade and telescopic power-driven arm. The chassis had suspension of M1A2 standard and was operated by a crew of two. Two vehicles were delivered to the Army in 1999. Development was halted by the following year. M1074 Joint Assault Bridge (JAB): Bridgelayer combining a heavy "scissor" bridge with the M1 Abrams chassis. Expected to reach low-rate initial production in 2019 to replace the M60 AVLB and M104 Wolverine. M1150 Assault Breacher Vehicle (ABV): Assault variant for the USMC. Based upon the M1A1 Abrams chassis, the Assault Breacher Vehicle has a variety of systems installed, such as a full-width mine plow, two linear demolition charges, and a lane-marking system. Reactive armor has been fitted to the vehicle providing additional protection against HEAT warhead-based weapons. The turret has been replaced by a new smaller one with two MICLIC launchers at its rear. A M2HB .50 machine gun in a remote weapons station is mounted on the commander's cupola and a bank of grenade launchers are fitted to each side of the superstructure to cover the frontal arc for self-protection. Additional equipment Mine clearing plows: An early example consisted of two mine plows and a chain with a weight running between them. The Mine Clearing Blade System (MCBS): It is capable of clearing mines up to 6 feet in front of the tank's path. Mine-clearing rollers: The Tank Mounted Mine Clearing Roller (TMMCR) comprises two roller banks of five discs each and a chain with a weight running between the rollers. Self Protection Combat Roller (SPCR) targets pressure activated explosive devices. The system comprises two 4-wheel roller gangs. An optional Magnetic System Duplicator (MSD) can be fitted to help protect the equipment from the effect of magnetic influence fused mines. Surface Clearance Device (SCD): The SCD is employed to clear surface laid mines and IEDs from roads, trails and rough terrain. There are two versions of the SCD; a V-blade optimised for clearing routes and a straight angle-blade which is optimised for clearing staging and assembly areas. Vehicle Magnetic Signature Duplicator (VEMSID): The VEMSID causes detonation of magnetic influence mines. The system comprises four emitter coils, two associated power boxes and an MSD Control Unit (MSDCU). Bulldozer attachment. The U.S. Army tested this attachment in 1982. This was unsuccessful in part because it resulted in transmission overheating, and was never used. Specifications Operators Current operators – Australian Army: 59 M1A1 as of 2024. AIM configuration tanks (hybrids with a mix of equipment used by U.S. Army and U.S. Marine Corps but without depleted uranium layers in armor). They were bought from the U.S. in 2006 and replaced the Leopard AS1 in 2007. As of 2017, the Australian Government was considering expanding the Army's fleet of Abrams to 90 tanks. In April 2021, the U.S. granted an FMS for 160 M1A1 tank hulls to produce 75 M1A2 SEPv3 tanks, 29 M1150 Assault Breacher Vehicles and 18 M1074 Joint Assault Bridges, including the development of a unique armor package for the Australian Army. In January 2022, Australia committed to purchase 120 tanks and armored vehicles, including 75 M1A2s, at a total cost of $3.5 billion and to be delivered in 2024; the M1A2s will replace their 59 M1A1s. On 17 October 2024, the Australian government announced that it will be gifting 49 of its retired M1A1 tanks to Ukraine. 14 M1A2s are scheduled to enter service by the end of 2024. – Egyptian Army: 1,130 M1A1 as of 2024. 1,360 M1A1 tanks assembled in Egypt for the Egyptian Army in cooperation with the U.S. – Iraqi Army: 100 M1A1 as of 2024. Iraq purchased 120 M1A1SA from the U.S. in 2008. The first 11 tanks were delivered to the Iraqi Army in August 2010 with all deliveries completed by August 2011. In October 2012, it was reported that six more tanks were being delivered. Four battalions of the 9th Armoured Division were equipped with M1s by 2014: 1st and 2nd of the 34th Brigade, and 4th and 5th of the 35th Brigade. Iraq purchased 175 more M1A1 in 2014, though it is unclear if these were delivered. – Kuwaiti Army: 218 M1A2K as of 2024. 218 M1A2s produced . Kuwait took delivery of the first of 217 M1A2K variants in 2021. – Royal Moroccan Army: 222 M1A1SA as of 2024. 222 M1A1 SA (situational awareness) tanks ordered in 2015. Deliveries under the contract started in July 2016 with an estimated completion date of February 2018. The contract include 150 refurbished and upgraded tanks to the special armor configuration. Morocco took delivery of the first batch of M1A1SAs on 28 July 2016. A Foreign Military Sale for 162 M1A2Ms was approved by the U.S. Department of State in November 2018 and sent to Congress for final approval. In October 2020 General Dynamics Land Systems was awarded a $11.9 million contract to upgrade 162 Abrams tanks to the M1A2 SEPv3 configuration. The contract was completed in March 2022. – Polish Land Forces: Poland purchased 116 former U.S. Marine Corps M1A1s in January 2023. The first 14 arrived on 28 June 2023. A further 26 were delivered in November 2023. On 8 January 2024 the next 29 were delivered. Deliveries were completed in June 2024. Poland also purchased 250 M1A2 SEPv3 tanks. Production is set to finish by 2024, and delivery in early 2025. 28 M1A2 SEPv2 tanks were leased in July 2022 to train crews until deliveries begin. – Saudi Arabian Army: 575 M1A2S as of 2024. 373 Abrams tanks first ordered to be upgraded to M1A2S configuration in Saudi Arabia. 69 more M1A2S tanks ordered on 8 January 2013, to be delivered by 31 July 2014. – Republic of China Army: In July 2018, Taiwan's Ministry of National Defense budgeted money to buy 108 M1A2Ts from the U.S. to replace its aging CM-11 Brave Tiger and M60A3 TTS tanks. The U.S. Department of State approved the $2.2 billion sale in July 2019. A sale of 108 M1A2Ts was later finalized. The first two were delivered to Taiwan in June 2022, with an additional 38 delivered in December 2024. – Ukrainian Ground Forces: 31 M1A1SA (Situational awareness). 31 M1A1 Abrams were delivered prior to 16 October 2023, as part of U.S. support for Ukraine. As of 21 March 2023, the US government is offering to supply Ukraine with older M1A1 Abrams tanks, "that have been upgraded very similar capability to the M1A2", as opposed to newer M1A2 tanks in an effort to speed up delivery. According to Pentagon Press Secretary U.S. Air Force Brigadier General Pat Ryder these tanks are to be delivered "by Fall" 2023 instead of the original delivery time of mid-2024. These tanks are to be taken from existing M1A1 "excessive hulls" from United States stocks and modernized to the required standard. In April, the U.S. announced that the Abrams tanks will soon be sent to Germany so that Ukrainian soldiers can start training on them. On 6 September, the U.S. said it would supply Ukraine with depleted uranium ammunition for Abrams tanks, despite initially refusing to do so. On 25 September, it was reported that Ukraine had received its first shipment of M1 tanks. On 16 October 2024, the Australian government announced a $245 million (AUD) aid package for Ukraine. It included the transfer of 49 M1A1 Abrams tanks to Ukraine. – United States Army received over 8,100 M1, M1A1 and M1A2 tanks combined. U.S. Army – 2,640 total in service – 540 M1A1 SA; 1,500 M1A2 SEPv2; 600 M1A2 SEPv3 s; (2,000 more M1A1/A2 Abrams in store). Future operators – Royal Bahraini Army: On 19 March 2024, the U.S. Defense Security Cooperation Agency announced that the Department of State had approved the possible Foreign Military Sale to Bahrain of 50 M1A2 SEPv3 Abrams Main Battle Tanks. – Romanian Land Forces: On 7 March 2023, a senior defense official announced that the Romanian Land Forces is in the process of advancing a proposal for the purchase of an Abrams tank battalion. In May 2023, the decision to buy 54 used M1A2 Abrams from United States Army stocks was approved by the Parliament of Romania. The Romanian M1A2 variant is designated M1A2R. On 9 November 2023, the U.S. Defense Security Cooperation Agency announced that the Department of State had approved the possible Foreign Military Sale to Romania for 54 M1A2 tanks and related equipment for an estimated cost of $2.53 billion. Former operators – United States Marine Corps: In 2020 the Marine Corps announced the disbandment of its tank units, citing a pivot towards amphibious warfare by implementing Force Design 2030. All 450 of the Marine Corps M1 Abrams MBTs were transferred to the U.S. Army with withdrawal from Marine Corps service being completed in May 2021. Former non-state operators – 9 vehicles were destroyed by ISIL Ramadi. – 1 vehicle was seen in use by Hezbollah around January 2015, another vehicle was again seen in use in January 2016 and another one was seen being transported in March 2016. PMF – 1 vehicle was seen in a PMF video montage.	Technology	Specific armor	null
37746	https://en.wikipedia.org/wiki/Boeing%20AH-64%20Apache	Boeing AH-64 Apache	The Boeing AH-64 Apache ( ) is an American twin-turboshaft attack helicopter with a tailwheel-type landing gear and a tandem cockpit for a crew of two. Nose-mounted sensors help acquire targets and provide night vision. It carries a M230 chain gun under its forward fuselage and four hardpoints on stub-wing pylons for armament and stores, typically AGM-114 Hellfire missiles and Hydra 70 rocket pods. Redundant systems help it survive combat damage. The Apache began as the Model 77 developed by Hughes Helicopters for the United States Army's Advanced Attack Helicopter program to replace the AH-1 Cobra. The prototype YAH-64 first flew on 30 September 1975. The U.S. Army selected the YAH-64 over the Bell YAH-63 in 1976, and later approved full production in 1982. After acquiring Hughes Helicopters in 1984, McDonnell Douglas continued AH-64 production and development. The helicopter was introduced to U.S. Army service in April 1986. The advanced AH-64D Apache Longbow was delivered to the Army in March 1997. Production has been continued by Boeing Defense, Space & Security. As of March 2024, over 5,000 Apaches have been delivered to the U.S. Army and 18 international partners and allies. Primarily operated by the U.S. Army, the AH-64 has also become the primary attack helicopter of multiple nations, including Greece, Japan, Israel, the Netherlands, Singapore, and the United Arab Emirates. It has been built under license in the United Kingdom as the AgustaWestland Apache. American AH-64s have served in conflicts in Panama, the Persian Gulf, Kosovo, Afghanistan, and Iraq. Israel uses the Apache to fight in Lebanon and the Gaza Strip. British and Dutch Apaches were deployed to wars in Afghanistan and Iraq in 2001. Development Advanced Attack Helicopter After the AH-56 Cheyenne was cancelled in 1972 in favor of projects like the U.S. Air Force A-10 Thunderbolt II and the Marine Corps AV-8A Harrier, the United States Army sought an aircraft to fill an anti-armor attack role that would still be under Army command. The 1948 Key West Agreement forbade the Army from owning combat fixed-wing aircraft. The Army wanted an aircraft better than the AH-1 Cobra in firepower, performance, and range. It would have the maneuverability for terrain following nap-of-the-earth (NoE) flying. To this end, the U.S. Army issued a Request For Proposals (RFP) for the Advanced Attack Helicopter (AAH) program on 15 November 1972. As a sign of the importance of this project, in September 1973 the Army designated its five most important projects as the "Big Five", with the AAH included. Proposals were submitted by Bell, Boeing Vertol/Grumman team, Hughes, Lockheed, and Sikorsky. In July 1973, the U.S. Department of Defense selected finalists Bell and Hughes Aircraft's Toolco Aircraft Division (later Hughes Helicopters). This began the phase 1 of the competition. Each company built prototype helicopters and went through a flight test program. Hughes' Model 77/YAH-64A prototype first flew on 30 September 1975, while Bell's Model 409/YAH-63A prototype first flew on 1 October 1975. After evaluating the test results, the Army selected Hughes' YAH-64A over Bell's YAH-63A in 1976. Reasons for selecting the YAH-64A included its more damage tolerant four-blade main rotor and the instability of the YAH-63's tricycle landing gear arrangement. The AH-64A then entered phase 2 of the AAH program under which three pre-production AH-64s would be built, additionally, the two YAH-64A flight prototypes and the ground test unit were upgraded to the same standard. Weapons and sensor systems were integrated and tested during this time, including the laser-guided AGM-114 Hellfire missile. Development of the Hellfire missile had begun in 1974, originally known by the name of Helicopter Launched, Fire and Forget Missile ('Hellfire' being a shortened acronym), for the purpose of arming helicopter platforms with an effective anti-tank missile. Into production In 1981, three pre-production AH-64As were handed over to the U.S. Army for Operational Test II. The Army testing was successful, but afterward it was decided to upgrade to the more powerful T700-GE-701 version of engine, rated at . The AH-64 was named the Apache in late 1981, after the Apache tribe, following the tradition of naming Army helicopters after Native American tribes. It was approved for full-scale production in 1982. In 1983, the first production helicopter was rolled out at Hughes Helicopter's facility at Mesa, Arizona. Hughes Helicopters was purchased by McDonnell Douglas for $470 million in 1984 (). The helicopter unit later became part of The Boeing Company with the merger of Boeing and McDonnell Douglas in August 1997. In 1986, the incremental or flyaway cost for the AH-64A was $7M and the average unit cost was approximately $13.9M based on total costs. A 1985 Department of Defense engineering analysis by the inspector general's office reported that significant design deficiencies still needed to be addressed by the contractor. The Army project manager Col. William H. Forster published a list of 101 action items. In 1986, the four 22-foot-long main rotor blades, each made from steel and composite material glued together to maximize strength and minimize weight by the Composite Structures Division of Alcoa Composites, were added to the list. The steel-composite rotors could not meet the Army specification for a life of 1500 flight hours, and needed replacement after just 146 hours. After six changes to the design, the rotor blade life was extended to 1400 hours by early 1991. As of 2024 the AH-64E is being produced at an economical rate of 82 aircraft a year. Boeing states that the minimum sustainment rate for the aircraft is 48 per year while current tooling and space allows for up to 98 aircraft to be manufactured per year. The U.S. Army states that with additional investment and labor production could be raised to 144 aircraft per year. Further development During the 1980s, McDonnell Douglas studied an AH-64B, featuring an updated cockpit, new fire control system and other upgrades. In 1988, funding was approved for a multi-stage upgrade program to improve sensor and weapon systems. Technological advance led to the program's cancellation in favor of more ambitious changes. In August 1990, development of the AH-64D Apache Longbow was approved by the Defense Acquisition Board. The first AH-64D prototype flew on 15 April 1992. Prototype testing ended in April 1995. During testing, six AH-64D helicopters were pitted against a bigger group of AH-64As. The results demonstrated the AH-64D to have a sevenfold increase in survivability and fourfold increase in lethality compared to the AH-64A. On 13 October 1995, full-scale production was approved; a $1.9-billion five-year contract was signed in August 1996 to upgrade 232 AH-64As into AH-64Ds. On 17 March 1997, the first production AH-64D flew. It was delivered on 31 March. Portions of the Apache are produced by other aerospace firms. AgustaWestland has produced a number of components for the Apache both for the international market and for the British Army's AgustaWestland Apache. Since 2004, Korea Aerospace Industries has been the sole manufacturer of the Apache's fuselage. Fuselage production had previously been performed by Teledyne Ryan Aeronautical. The transfer of fuselage production led to a prolonged legal dispute between Teledyne Ryan and Boeing. The AH-64D program cost a total of $11 billion (~$ in ) through 2007. In April 2006, Boeing was awarded a $67.6 million (~$ in ) fixed-price contract for the remanufacture of several existing U.S. AH-64As to the AH-64D configuration. Between May 2009 and July 2011, a further five contracts were issued to remanufacture batches of AH-64As into AH-64Ds. Since 2008, nations operating the older Apaches have been urged to undertake modernization programs as support for the AH-64A is withdrawn. By May 2019, Boeing tested in a wind tunnel a compound Apache scale model with a pusher propeller, a small wing to increase range and speed, and a counter-torque tail rotor like the cancelled Lockheed AH-56 Cheyenne of the 1960s. It competed for the U.S. Army's FLRAA unveiled in April, developed from the Future Vertical Lift Capability Set 3 (medium rotorcraft) without the attack requirement, while the U.S. Army's FARA should replace the retired Bell OH-58 Kiowa scout and up to half of the AH-64 fleet. Design Overview The AH-64 Apache has a four-blade main rotor and a four-blade tail rotor. The crew sits in tandem, with the pilot sitting behind and above the co-pilot/gunner. Both crew members are capable of flying the aircraft and performing methods of weapon engagements independently. The AH-64 is powered by two General Electric T700 turboshaft engines with high-mounted exhausts on either side of the fuselage. Various models of engines have been used on the Apache; those in British service use engines from Rolls-Royce. In 2004, General Electric Aviation began producing more powerful T700-GE-701D engines, rated at for AH-64Ds. The crew compartment and rotor blades are designed to sustain a hit from rounds. The airframe includes some of protection and has a self-sealing fuel system to protect against ballistic projectiles. The crew compartment also incorporates a transparent blast shield between the pilot and gunner seats so that at least one crew member can survive in the event of a direct hit, but the canopy and windows are otherwise unrated against ballistic threats. The aircraft was designed to meet the crashworthiness requirements of MIL-STD-1290, which specifies minimum requirement for crash impact energy attenuation to minimize crew injuries and fatalities. This was achieved through incorporation of increased structural strength, crashworthy landing gear, seats and fuel system. On a standard day, when temperatures are 59 °F (15 °C), the AH-64 has a vertical rate of climb of 1,775 feet per minute (541 m/min), and a service ceiling of 21,000 feet (6,400 m). However, on a hot day, when temperatures are 70 °F (21 °C), its vertical rate of climb is reduced to 1,595 fpm (486 m/min), and service ceiling is reduced to 19,400 feet (5,900 m) due to less dense air. Avionics and targeting One of the revolutionary features of the Apache was its helmet mounted display, the Integrated Helmet and Display Sighting System (IHADSS); among its capabilities, either the pilot or gunner can slave the helicopter's 30 mm automatic M230 Chain Gun to their helmet, making the gun track head movements to point where they look. The M230E1 can be alternatively fixed to a locked forward firing position, or controlled via the Target Acquisition and Designation System (TADS). On more modern AH-64s, the TADS/PNVS has been replaced by Lockheed Martin's Arrowhead (MTADS) targeting system. U.S. Army engagement training is performed under the Aerial Weapons Scoring System Integration with Longbow Apache Tactical Engagement Simulation System (AWSS-LBA TESS), using live 30 mm and rocket ammunition as well as simulated Hellfire missiles. The Smart Onboard Data Interface Module (SMODIM) transmits Apache data to an AWSS ground station for gunnery evaluation. The AH-64's standard of performance for aerial gunnery is to achieve at least 1 hit for every 30 shots fired at a wheeled vehicle at a range of . The AH-64 was designed to perform in front-line environments, and to operate at night or day and during adverse weather conditions, thanks to systems including the Target Acquisition and Designation System, Pilot Night Vision System (TADS/PNVS), passive infrared countermeasures, GPS, and the IHADSS. Longbow-equipped Apaches can locate up to 256 targets simultaneously within . In August 2012, 24 U.S. Army AH-64Ds were equipped with the Ground Fire Acquisition System (GFAS), which detects and targets ground-based weapons fire sources in all-light conditions and with a 120° visual field. The GFAS consists of two sensor pods working with the AH-64's other sensors, and a thermographic camera that precisely locates muzzle flashes. In 2014, it was announced that new targeting and surveillance sensors were under development to provide high-resolution color imagery to crews, replacing older low definition black-and-white imaging systems. Lockheed received the first contract in January 2016, upgrading the Arrowhead turret to provide higher-resolution color imaging with longer ranges and a wider field of view. In 2014, the U.S. Army was adapting its Apaches for increased maritime performance as part of the Pentagon's rebalance to the Pacific. Additional avionics and sensor improvements includes an extended-range radar capable of detecting small ships in littoral environments, software adaptions to handle maritime targets, and adding Link 16 data-links for better communications with friendly assets. Armament and configurations The AH-64 is adaptable to numerous different roles within its context as Close Combat Attack (CCA). In addition to the 30 mm M230E1 Chain Gun, the Apache carries a range of external stores and weapons on its stub-wing pylons, typically a mixture of AGM-114 Hellfire anti-tank missiles, and Hydra 70 general-purpose unguided rockets. The Hellfire is designed to defeat stationary or moving tanks as far away as 6,500 meters. Since 2005, the Hellfire missile outfitted with a thermobaric warhead is designated AGM-114N; this missile version is intended for use against ground forces and urban warfare operations. In October 2015, the U.S. Army ordered its first batch of Advanced Precision Kill Weapon System (APKWS) guided 70 mm rockets for the Apache. Starting in the 1980s, the Stinger and AIM-9 Sidewinder air-to-air missiles and the AGM-122 Sidearm anti-radiation missile were evaluated for use upon the AH-64. The Stinger was initially selected; the U.S. Army was also considering the Starstreak air-to-air missile. External fuel tanks can also be carried on the stub wings to increase range and mission time. The stub-wing pylons have mounting points for maintenance access; these mountings can also be used to secure personnel externally for emergency transport. Stinger missiles are often used on non-U.S. Apaches, as foreign forces do not have as many air superiority aircraft to control the skies. The AH-64E initially lacked the ability to use the Stinger to make room for self-defense equipment, but the capability was added back following a South Korean demand. The AH-64E is able to control unmanned aerial vehicles (UAVs), used by the U.S. Army to perform aerial scouting missions previously performed by the OH-58 Kiowa. Apaches can request to take control of an RQ-7 Shadow or MQ-1C Grey Eagle from ground control stations to safely scout via datalink communications. There are four levels of UAV interoperability (LOI): LOI 1 indirectly receives payload data; LOI 2 receives payload data through direct communication; LOI 3 deploys the UAV's armaments; and LOI 4 takes over flight control. UAVs can search for enemies and, if equipped with a laser designator, target them for the Apache or other friendly aircraft. Boeing has suggested that the AH-64 could be fitted with a directed energy weapon. The company has developed a small laser weapon, initially designed to engage small UAVs, that uses a high-resolution telescope to direct a 2–10 kW beam with the diameter of a penny out to a range of . On the Apache, the laser could be used to destroy enemy communications or radio equipment. On 26 June 2017, the Army and Raytheon announced they had successfully completed the first-ever helicopter-based flight demonstration of a high energy laser system from an AH-64. On 14 July 2016, it was reported that the AH-64 had successfully completed testing of the MBDA Brimstone anti-armor missile. In January 2020, the U.S. Army announced it was fielding the Spike NLOS missile on AH-64E Apaches as an interim solution to acquire new munitions that provide greater stand-off capabilities. Operational history United States Twentieth century In January 1984, the U.S. Army formally accepted its first production AH-64A and training of the first pilots began later that year. The first operational Apache unit, 7th Battalion, 17th Cavalry Brigade, began training on the AH-64A in April 1986 at Fort Hood, Texas. Two operational units with 68 AH-64s first deployed to Europe in September 1987 and took part in large military exercises there. Upon fielding the Apache, capabilities such as the FLIR's use in extensive night operations made it clear that it was capable of operating beyond the forward line of own troops (FLOT) to which previous attack helicopters were normally restricted. It was discovered that the Apache was coincidentally fitted with the Have Quick UHF radio system used by the U.S. Air Force, allowing inter-service coordination and joint operations such as the joint air attack teams (JAAT). The Apache has operated extensively with close air support (CAS) aircraft, such as the USAF's Fairchild Republic A-10 Thunderbolt II and the USMC's McDonnell Douglas AV-8B Harrier II, often acting as a target designator to conserve the Apache's own munitions. The Apache was first used in combat in 1989, during Operation Just Cause, the invasion of Panama. It participated in over 240 combat hours, attacking various targets, mostly at night. General Carl Stiner, the commander of the operation, stated: "You could fire that Hellfire missile through a window from four miles away at night." Nearly half of all U.S. Apaches were deployed to Saudi Arabia following Iraq's invasion of Kuwait in 1990. During Operation Desert Storm on 17 January 1991, eight AH-64As guided by four MH-53 Pave Low IIIs destroyed part of Iraq's radar network in the operation's first attack, allowing the attack aircraft to evade detection. Each Apache carried an asymmetric load of Hydra 70 rockets, Hellfires, and one auxiliary fuel tank. During the 100-hour ground war, a total of 277 AH-64s took part, destroying 278 tanks, numerous armored personnel carriers and other Iraqi vehicles, for a total of over 500 kills. One AH-64 was lost in the war, crashing after a close range rocket-propelled grenade (RPG) hit; the crew survived. While effective in combat, the AH-64 presented serious logistical difficulties. Findings reported in 1990 stated "maintenance units could not keep up with the Apache's unexpectedly high work load..." To provide spare parts for combat operations, the U.S. Army unofficially grounded all other AH-64s worldwide; Apaches in the theater flew only one-fifth of planned flight-hours. Such problems were evident before the Gulf War. The AH-64 played roles in the Balkans during separate conflicts in Bosnia and Kosovo in the 1990s. During Task Force Hawk, 24 Apaches were deployed to a land base in Albania in 1999 for combat in Kosovo. These required 26,000 tons of equipment to be transported over 550 C-17 flights, at a cost of . During these deployments, the AH-64 encountered problems, such as deficiencies in training, night vision equipment, fuel tanks, and survivability. In 2000, Major General Dick Cody, 101st Airborne's commanding officer, wrote a strongly worded memo to the Chief of Staff about training and equipment failures. Almost no pilots were qualified to fly with night vision goggles, preventing nighttime operations. The Washington Post printed a front-page article on the failures, commenting: "The vaunted helicopters came to symbolise everything wrong with the Army as it enters the 21st century: Its inability to move quickly, its resistance to change, its obsession with casualties, its post-Cold War identity crisis". Although no Apache combat missions took place, two were lost in training exercises. An effective network of Yugoslav air defenses stopped Apaches from being deployed on combat missions in Kosovo. 21st century U.S. Apaches served in Operation Enduring Freedom in Afghanistan from 2001. It was the only Army platform capable of providing accurate CAS duties for Operation Anaconda, often taking fire and quickly repaired during the intense early fighting. Apaches often flew in small teams with little autonomy to react to threats and opportunities, requiring lengthy dialogue with centrally micromanaged command structures. U.S. AH-64Ds typically flew in Afghanistan and Iraq without the Longbow Radar in the absence of armored threats. On 21 December 2009, a pair of U.S. Apaches attacked a British-held base in a friendly fire incident, killing one British soldier. In 2003, the AH-64 participated in the invasion of Iraq during Operation Iraqi Freedom. On 24 March 2003, 31 Apaches were damaged; one was shot down in an unsuccessful attack on an Iraqi Republican Guard armored brigade near Karbala. Iraqi tank crews had set up a "flak trap" among terrain and effectively employed their guns. Iraqi officials claimed a farmer with a Brno rifle shot down the Apache, but the farmer denied involvement. The AH-64 came down intact and the crew were captured; it was destroyed via air strike the following day. This incident had significant consequences for the AH-64 helicopter because it revealed an important vulnerability. Despite being considered by army aviators as flying tanks at the time, it became clear that the AH-64 was actually highly susceptible to rifle fire. As a result, the army quietly disclosed in early 2006 that AH-64s would no longer have a major role in carrying out attacks deep inside enemy lines. By the end of U.S. military operations in Iraq in December 2011, several Apaches had been shot down by enemy fire and lost in accidents. In 2006, an Apache was downed by a Soviet-made Strela 2 (SA-7) in Iraq, despite it being typically able to avoid such missiles. In 2007, four Apaches were destroyed on the ground by insurgent mortar fire using web-published geotagged photographs taken by soldiers. Several AH-64s were lost to accidents in Afghanistan. Most Apaches that took heavy damage were able to continue their missions and return safely. By 2011, the U.S. Army Apache fleet had accumulated more than 3 million flight hours since the first prototype flew in 1975. A DOD audit released in May 2011 found that Boeing had frequently overcharged the U.S. Army for routine spare parts in helicopters like the Apache, ranging from 33.3 percent to 177,475 percent. On 21 February 2013, the 1st Battalion (Attack), 229th Aviation Regiment at Joint Base Lewis–McChord became the first U.S. Army unit to field the AH-64E Apache Guardian; a total of 24 AH-64E were received by mid-2013. On 27 November 2013, the AH-64E achieved initial operating capability (IOC). In March 2014, the 1st–229th Attack Reconnaissance Battalion deployed 24 AH-64Es to Afghanistan in the type's first combat deployment. From April through September 2014, AH-64Es in combat maintained an 88 percent readiness rate. The unit's deployment ended in November 2014, with the AH-64E accumulating 11,000 flight hours, each helicopter averaging 66 hours per month. The AH-64E flies faster than the AH-64D, cutting response time by 57 percent, and has better fuel efficiency, increasing time on station from 2.5–3 hours to 3–3.5 hours; Taliban forces were reportedly surprised by the AH-64E attacking sooner and for longer periods. AH-64Es also worked with medium and large unmanned aerial vehicles (UAVs) to find targets and maintain positive ID, conducting 60 percent of the unit's direct-fire engagements in conjunction with UAVs; Guardian pilots often controlled UAVs and accessed their video feeds to use their greater altitudes and endurance to see the battlespace from standoff ranges. In 2014, the Army began implementing a plan to move all Apaches from the Army Reserve and National Guard to the active Army to serve as scout helicopters to replace the OH-58 Kiowa. Using the AH-64 to scout would be less expensive than Kiowa upgrades or purchasing a new scout helicopter. AH-64Es can control UAVs like the MQ-1C Grey Eagle to perform aerial scouting missions; a 2010 study found the teaming of Apaches and UAVs was the most cost-effective alternative to a new helicopter and would meet 80 percent of reconnaissance requirements, compared to 20 percent with existing OH-58s and 50 percent with upgraded OH-58s. National Guard units, who would lose their attack helicopters, criticized the proposal. In March 2015, the first heavy attack reconnaissance unit was formed with 24 Apaches and 12 Shadow UAVs. In July 2014, the Pentagon announced that Apaches had been dispatched to Baghdad to protect embassy personnel from Islamic State militant attacks. On 4 October 2014, Apaches began performing missions in Operation Inherent Resolve against Islamic State ground forces. In October 2014, U.S. Army AH-64s and Air Force fighters participated in four air strikes on Islamic State units northeast of Fallujah. In June 2016, Apaches were used in support of the Iraqi Army's Mosul offensive and provided support during the Battle of Mosul, sometimes flying night missions supporting Iraqi operations. In December 2019, two Apaches provided overwatch for U.S. Marines to secure the U.S. embassy in Baghdad, Iraq after armed militants, supported by Iran, attempted to storm the facility. In March 2024, two Apache crashes within two days increased scrutiny and made national news. Israel The Israeli Air Force (IAF) first received AH-64As in 1990, for a fleet of 42 by 2000. The IAF's choice to buy Apaches over upgrading its AH-1 Cobra attack helicopters was controversial. In 2000, Israel was interested in acquiring up to 48 AH-64Ds, but U.S. reluctance to share the source code complicated the prospect. In April 2005, Boeing delivered the IAF's first AH-64D. In 2001, the U.S. government was allegedly investigating misuse of the Apache and other U.S.-supplied military equipment against Palestinians. In 2009, the sale of six AH-64Ds was reportedly blocked by the Obama Administration, pending interagency review, over concerns they may pose a threat to civilian Palestinians in Gaza. In IAF service, the AH-64A was named Peten (, for Cobra), while the AH-64D was named Saraph (, for venomous/fiery winged serpent). During the 1990s, Israeli AH-64As frequently attacked Hezbollah outposts in Lebanon. On 13 April 1996, during Operation Grapes of Wrath, an Apache fired two Hellfire missiles at an ambulance in Lebanon, killing six civilians. During the al-Aqsa Intifada in the 2000s, AH-64s were used to kill senior Hamas figures, such as Ahmed Yassin, Abdel Aziz al-Rantisi, and Adnan al-Ghoul. Human Rights Watch documented instances of IAF Apaches attacking civilian homes during the 2002 Jenin operation, killing one civilian. Consequently, HRW urged the US government to seek written assurances from Israel that Apaches would not be used to violate humanitarian law in the future. In 2004, Israeli AH-64s carried out the assassination of Ahmed Yassin and also killed 7 bystanders. Ahmed Yassin was the spiritual leader of Hamas; given that he was also blind, paraplegic and in a wheelchair, Palestinians saw the killing as "a cowardly execution of a frail old man in a wheelchair who did not attempt to hide". The attack also killed 7 bystanders, and was internationally condemned. IAF Apaches played a prominent role in the 2006 Lebanon War, launching strikes into Lebanon targeting Hezbollah forces. IAF Apaches also attacked civilian targets, killing many, including women and children. During this war, two AH-64As collided, killing one pilot and critically wounding three. In another incident in the conflict, an IAF AH-64D crashed due to a main rotor malfunction, killing the two crew. In late 2007, the IAF put further purchases and deliveries of AH-64Ds on hold while its performance envelope was investigated. Israeli officials praised the Apache for its role in Operation Cast Lead in 2008, against Hamas in Gaza. IAF Apaches have often patrolled the skies over Gaza; strikes against insurgents by these helicopters has become a frequent occurrence. In the 2010s, the IAF pursued upgrades to its AH-64A fleet as new AH-64D orders had been blocked. In June 2010, Israel decided not to upgrade all AH-64As to AH-64Ds due to funding constraints and lack of U.S. cooperation. In December 2010, the IAF was examining the adoption of a new missile system as a cheaper and lightweight complement to the Hellfire missile, either the American Hydra 70 or the Canadian CRV7. By 2013, IAF AH-64As were receiving a comprehensive upgrade of their avionics and electrical systems. The AH-64As are being upgraded to the AH-64Ai configuration, which is near the AH-64D standard. IAF Apaches can carry Spike anti-tank missiles instead of the Hellfire. The latest AH-64D-I integrates Israeli systems such as Elta communications suite, Elbit mission management system, Rafael Combat Net system and Elisra self-protection suite. IAF AH-64s occasionally saw use in the air-to-air role. The first operational air-to-air kill took place on 24 May 2001, when an IAF shot down a Lebanese civilian Cessna 152 aircraft. Israeli and Lebanese officials presented differing versions: Lebanon said Israel first intercepted the aircraft over Lebanese airspace and its pilot, flying without his instructor, mistakenly entered Israeli airspace, while Israel says the aircraft was already in Israeli airspace when it was intercepted and repeatedly refused to answer or comply with air traffic control (ATC) warnings. The second air-to-air kill occurred on 10 February 2018, after an Iranian UAV entered Israeli airspace from Syria, an AH-64 destroyed it with a missile. United Kingdom The UK previously operated a modified version of the AH-64D Block I Apache Longbow; initially called the Westland WAH-64 Apache, it is designated the Apache AH1 by the British Army. Westland built 67 WAH-64 Apaches under license from Boeing, following a competition between the Eurocopter Tiger and the Apache for the British Army's new Attack Helicopter in 1995. Important deviations made by AgustaWestland from the U.S. Apache variants include changing to more powerful Rolls-Royce engines, and the addition of a folding blade assembly for use on naval ships. On 11 July 2016, the Ministry of Defence confirmed a U.S. Foreign Military Sale (FMS) worth $2.3 billion (~$ in ) for 50 AH-64Es to be built in Mesa, Arizona. Leonardo Helicopters in the UK will maintain the current fleet of Apaches until 2023–2024, with a long-term plan for Leonardo and other UK companies to "do most of the work" on the new fleet. The deal included an initial support contract for maintenance, spare parts, and training simulators; components from the older WAH-64s "will be reused and incorporated into the new helicopters where possible." The type entered service with the British Army in 2022. Approval for the re-manufacture of fifty of the UK's WAH-64 Mk 1 fleet to AH-64E Apache Guardian standard was given by the Defense Security Cooperation Agency in August 2015. They utilize the General Electric T700 engine rather than the Turbomeca RTM322 of the WAH-64; the first off-the-shelf purchase of a GE engine by the Ministry of Defence. The first two AH-64Es were delivered to the British Army on 26 November 2020. The older AH1 (WAH-64) were retired by 2024 in favour of the AH-64E models. Netherlands The Dutch government initially showed an interest in acquiring Apache helicopters in the late 1980s, when it stated that it may purchase as many as 52. A competition held in 1994 against the Eurocopter Tiger and the Bell AH-1 SuperCobra led to the Royal Netherlands Air Force ordering 30 AH-64D Apaches in 1995. Deliveries began in 1998 and ended in 2002. The RNLAF Apaches are equipped with the Apache Modular Aircraft Survivability Equipment (AMASE) self-protection system to counter infrared (IR) missile threats. The RNLAF Apaches' first deployment was in 2001 to Djibouti, Africa. They were also deployed alongside U.S. AH-64s in support of NATO peacekeeping forces in Bosnia and Herzegovina. In 2004, six Dutch AH-64s were deployed as part of the Netherlands contribution to the Multinational force in Iraq to support the Dutch ground forces. The Apaches performed close combat support and display of force missions, along with providing reconnaissance information to ground forces. In February 2006, the Netherlands' contribution to NATO forces in Afghanistan was increased from 600 to 1,400 troops and 6 AH-64s were sent in support. Shortly after Apaches were deployed to Hamid Karzai International Airport, as part of the Dutch contribution to ISAF, on 10 April 2004, a pair of Dutch Apaches came under light gunfire close to Kabul. On 17 December 2007, an RNLAF Apache flew into power lines during a night flying exercise in the Netherlands, forcing an emergency landing and causing a lengthy blackout in the region. On 17 March 2015, a RNLAF Apache crashed during a training mission in Mali. Both pilots died. The Ministry of Defence opened an investigation into the cause of the crash. In February 2018, the Netherlands decided to upgrade all their AH-64Ds to the latest AH-64E standard via a FMS contract with the US, along with 17 APG-78 fire control radar units. In November 2021, the process to upgrade AH-64Ds began and RNLAF is to receive the upgraded AH-64Es between 2023 and 2025. Saudi Arabia Following the 1991 Gulf War, during which many U.S. Apaches operated from bases within Saudi territory, Saudi Arabia purchased twelve AH-64As for the Royal Saudi Land Force. It has been speculated that the Saudi purchase motivated Israel to also procure Apaches. In August 2006, Saudi Arabia began negotiations for Apache upgrades worth up to $400M (~$ in ), possibly remanufacturing their AH-64As to the AH-64D configuration. In September 2008, the U.S. Government approved a Saudi Arabian request to buy 12 AH-64Ds. In October 2010, Saudi Arabia requested a further 70 AH-64Ds as part of a possible massive arms deal. In November 2009, the Royal Saudi Land Force, as part of a military effort against insurgent border intrusions, launched Operation Scorched Earth; this involved Apaches launching air strikes against Houthi rebels operating inside neighboring Yemen. In January 2010 the rebels claimed to have shot down an Apache; this was denied by the Saudi military. In late January 2010, the leader of the Shiite rebels announced their withdrawal from Saudi territory, this announcement followed a key battle on 12 January when Saudi forces reportedly took control of the border village of Al Jabiri. As an escalation of the Yemeni Civil War, starting on 26 March 2015, Saudi Arabia, the United Arab Emirates, and several other regional allies started a military operation in Yemen. Both the Saudi Army Aviation and the United Arab Emirates Air Force used their AH-64s in combat against an alliance between elements of the Yemeni Army loyal to the former president Saleh and the Houthis. The Apaches were mostly involved in border patrol and strikes in Northwestern Yemen. Over the years, several Saudi and an Emirati Apaches were lost to incidents and enemy fire, although exact numbers have not been independently confirmed. On 17 March 2017, an Apache reportedly attacked a Somali refugee boat, killing 42 refugees. Saudi Arabia denied involvement, though it and the United Arab Emirates are the only militaries using Apaches during the conflict. United Arab Emirates The United Arab Emirates purchased 30 AH-64As between 1991 and 1994, and began upgrading to AH-64D specification in 2008. In December 2016, the U.S. State Department approved a proposed sale of another 37 AH-64Es and Congress was notified; this consisted of 28 re-manufactured and nine new-build helicopters. A UAE AH-64 was reportedly lost on 17 October 2017; a replacement was approved by the US in 2019. Egypt In 1995, the Egyptian Air Force placed an order for 36 AH-64As. These Apaches were delivered with the same avionics as the U.S. fleet at that time, except for indigenous radio equipment. In 2000, Boeing announced an order to remanufacture Egypt's Apache fleet to the AH-64D configuration, except for Longbow radar, which had been excluded by the U.S. government. In 2009, Egypt requested a further 12 AH-64D Block IIs with Longbow radars through a FMS. In August 2012, the Egyptian Armed Forces undertook a large-scale military operation to regain control of the Sinai Peninsula from armed militants. Air cover throughout the operation was provided by the Egyptian Air Force's Apaches, which reportedly destroyed three vehicles and killed at least 20 militants. Up to five Egyptian Apaches were temporarily stationed in the Sinai following an agreement between Egypt and Israel. In September 2015, an Egyptian Apache attacked a group of foreign tourists in the Western Desert, killing 12 people and injuring 10. The AH-64s fired at the civilians with rockets and 30mm machine guns for several hours, even though survivors said they waved the white flag. The Egyptian Interior Ministry stated that the group, whom were mistaken for militants, were in a restricted area. The tourists were reportedly accompanied by Egyptian police, and their vehicles were marked with logos of the tourist company. In November 2018, the U.S. Department of State approved the sale of ten AH-64Es and associated equipment to Egypt, valued at US$1 billion, pending the sale clearing Congress. The Apache entered service with Egyptian Armed Services, and in the 2020s was being further upgraded. India Indian Air Force In 2008, the Indian Air Force (IAF) released a tender for 22 attack helicopters; there were six contending submissions: Sikorsky's UH-60 Black Hawk, the AH-64D, Bell's AH-1 Super Cobra, Eurocopter's Tiger, Mil's Mi-28 and AgustaWestland's A129 Mangusta. In October 2008, Boeing and Bell withdrew. In 2009, the competition was restarted. In December 2010, India requested the sale of 22 Apaches and associated equipment. On 5 October 2012, IAF Chief NAK Browne confirmed the Apache's selection. The IAF sought control of the 22 Apaches for air combat missions, while the Army Aviation Corps argued that they would be better used in army operations. In April 2013, the Indian Ministry of Defence decided that the 22 AH-64s would go to the IAF. India ordered the 22 AH-64Es in 2015. On 11 May 2019, the IAF received its first AH-64E in a ceremony at Boeing's Mesa, Arizona facility. On 3 September 2019, 8 AH-64Es were inducted into the IAF's 125 Helicopter Squadron at Pathankot Air Base, Punjab. Indian Army On 12 June 2018, the U.S. Department of State approved a possible Foreign Military Sales (FMS) to India for six more AH-64Es and associated equipment in an estimated $930 million deal. The U.S. Defense Security Cooperation Agency notified Congress for approval. In February 2020, another six for the Indian Army were ordered, including weapons, equipment, and training; deliveries are planned to begin in 2023. These attack helicopters are often interlinked with squadrons of indigenous HAL Prachand attack helicopters. On 1 January 2024, senior Army officials told India Today that the Indian Army Aviation Corps is expected to induct the first batch of Apache helicopters in February–March of this year. This is to enable the army to protect its tanks on the battlefield when the Indian Air Force is unavailable. They will be deployed to Jodhpur, near the India–Pakistan border, enhancing the security of the area against Pakistani tanks. On 15 March 2024, Army Aviation Corps raised an 451 Army Aviation squadron at Jodhpur which will operate Apache. The induction of first batch of three Apache was scheduled in May 2024 and the rest by July 2024. As of August 2024, no Apaches were delivered to the Army. Citing delays due to supply chain issues, as of September 2024, the first batch of three Apache helicopters are to delivered by the following December–January followed by the next 3 within another few months. Other users Greece received 20 AH-64As by 1995; another 12 AH-64Ds were ordered in September 2003. Singapore purchased 20 AH-64Ds in two batches between 1999 and 2001. In October 2010, training was suspended following the forced crash-landing of an Apache. In 2005, Kuwait purchased 16 AH-64Ds. On 26 August 2013, Indonesia and the U.S. finalized a contract for eight AH-64E Apaches worth $500 million (~$ in ). The first was displayed on 5 October 2017 as part of a military exercise in Indonesia, to mark the 72nd anniversary of its armed forces. The first batch of AH-64s for the Indonesian Army arrived in Indonesia on 18 December 2017. Japan ordered 50 AH-64Ds, which were built under license by Fuji Heavy Industries, designated "AH-64DJP". The first helicopter was delivered to the JGSDF in early 2006. The order was halted after 13 aircraft were delivered due to cost. In 2017, it was announced that the targeting systems of the 13 aircraft would be upgraded. One was destroyed in a crash in February 2018 with the deaths of both crew. In June 2011, Taiwan (Republic of China) agreed to the purchase of 30 AH-64Es with weapons and associated equipment. On 5 November 2013, Taiwan received the first six AH-64s. On 25 April 2014, a Taiwanese AH-64E crashed into a three-story building during a training flight in bad weather conditions, the first AH-64E airframe loss. An investigation ruled out mechanical failure and concluded that the pilots had rapidly descended through clouds at low altitude without checking flight instruments to maintain adequate height; the Army expanded simulator training in response. In October 2014, Taiwan's fifth and final batch of AH-64Es was delivered. In 2009, South Korea showed interest in the Apache, potentially related to plans to withdraw many U.S. Apaches from the country. On 21 September 2012, the U.S. Congress was notified of the possible purchase of 36 AH-64E, along with associated equipment and armament. It competed against the Bell AH-1Z Viper and the TAI/AgustaWestland T-129; in April 2013, South Korea announced plans to buy 36 AH-64Es. The first four AH-64Es were delivered in May 2016, and all 36 were deployed by January 2017. Future and possible users Iraq requested the sale of 24 AH-64s in April 2013; In January 2014, a sale, including the helicopters, associated parts, maintenance, and training, was cleared by the U.S. Congress. However, the proposal was not accepted by the Iraqi government and expired in August 2014. In July 2012, Qatar requested the sale of 24 AH-64Es with associated equipment and support. The sale was approved on 27 March 2014. In March 2019, Qatar received the first of 24 AH-64Es ordered. In July 2016, the UK placed an order for 50 AH-64Es through the U.S. FMS program instead of upgrading their Westland-built WAH-64s. In July 2019, Australia issued a request for information for Project Land 4503 to replace the Army's Eurocopter Tiger ARH helicopters. On 15 January 2021, the Australian Minister for Defence Linda Reynolds announced that the AH-64E had been selected to replace the Tiger ARH. A fleet of 29 AH-64Es will be acquired with a planned initial operational capability of 12 helicopters in 2026 and full operational capability in 2028. In October 2019, the Bangladesh Air Force (BAF) was offered two types of attack helicopters and selected the AH-64, pending government approval. However, due to the Apache's high costs, the BAF choose the competing Russian Mi-28NE Night Hunter. In November 2019, the U.S. State Department approved a FMS to Morocco of 24 AH-64Es (with an option for a further 12), this allows Morocco to negotiate an order. On 30 April 2020, the U.S. Defense Security Cooperation Agency announced it had received U.S. State Department approval and notified Congress of a possible sale to the Philippines of either six AH-1Z attack helicopters and related equipment for an estimated cost of $450 million or six AH-64Es and related equipment for an estimated cost of $1.5 billion. On 21 April 2022, Polish Ministry of National Defence shortlisted two models in competition for the "Kruk" (Raven) Program aimed at modernizing the Polish Land Forces' fleet of attack helicopters, with the competitors being the AH-64E and Bell's AH-1Z Viper. On 8 September 2022, Polish Minister of Defence Mariusz Błaszczak announced that the AH-64E has won and set out to procure 96 helicopters to form six squadrons. The contract includes a logistics package and a training package along with a stock of ammunition and spare parts. Variants AH-64A The AH-64A is the original production attack helicopter. The crew sit in tandem in an armored compartment. It is powered by two GE T700 turboshaft engines. The A-model was equipped with the −701 engine version until 1990 when the engines were switched to the more powerful −701C version. U.S. Army AH-64As are being converted to AH-64Ds. The service's last AH-64A was taken out of service in July 2012 before conversion at Boeing's facility in Mesa, Arizona. On 25 September 2012, Boeing received a $136.8M contract to remanufacture the last 16 AH-64As into the AH-64D Block II version and this was planned to be completed by December 2013. AH-64B In 1991, after Operation Desert Storm, the AH-64B was a proposed upgrade to 254 AH-64As. The upgrade would have included new rotor blades, a Global Positioning System (GPS), improved navigation systems and new radios. U.S. Congress approved $82M to begin the Apache B upgrade. The B program was canceled in 1992. The radio, navigation, and GPS modifications were later installed on most AH-64As via other upgrades. AH-64C Additional funding from Congress in late 1991 resulted in a program to upgrade AH-64As to an AH-64B+ version. More funding changed the plan to upgrade to AH-64C, which would include all of the changes to be included in the AH-64D except for mast-mounted radar and newer −700C engine versions. However, the C designation was dropped after 1993. With AH-64As receiving the newer engine from 1990, the only difference between the AH-64C and the AH-64D was the radar, which could be moved from one aircraft to another; thus, the decision was made to simply designate both versions as AH-64D. AH-64D The AH-64D Apache Longbow is equipped with a glass cockpit and advanced sensors, the most noticeable of which being the AN/APG-78 Longbow millimeter-wave fire-control radar (FCR) target acquisition system and the Radar Frequency Interferometer (RFI), housed in a dome located above the main rotor. The radome's raised position enables target detection while the helicopter is behind obstacles (e.g. terrain, trees or buildings). The AN/APG-78 is capable of simultaneously tracking up to 128 targets and engaging up to 16 at once; an attack can be initiated within 30 seconds. A radio modem integrated with the sensor suite allows data to be shared with ground units and other Apaches, allowing them to fire on targets detected by a single helicopter. The aircraft is powered by a pair of uprated T700-GE-701C engines. The forward fuselage was expanded to accommodate new systems to improve survivability, navigation, and 'tactical internet' communications capabilities. In February 2003, the first Block II Apache was delivered to the U.S. Army, featuring digital communications upgrades. The Japanese Apache AH-64DJP variant is based on the AH-64D; it can be equipped with the AIM-92 Stinger air-to-air missiles for self-defense. AH-64E Formerly known as AH-64D Block III, in 2012, it was redesignated as AH-64E Guardian. It has improved digital connectivity, the Joint Tactical Information Distribution System, more powerful T700-GE-701D engines with upgraded face gear transmission to handle more power, capability to control unmanned aerial vehicles (UAVs), full IFR capability, and improved landing gear. New composite rotor blades, which completed testing in 2004, increase cruise speed, climb rate, and payload capacity. Deliveries began in November 2011. Full-rate production was approved on 24 October 2012. The total Army Acquisition Objective for both new build and remanufactured AH-64Es is 812. Changes in production lots 4 through 6 include a cognitive decision aiding system and new self-diagnostic abilities. The updated Longbow radar has an oversea capacity, potentially enabling naval strikes; an AESA radar is under consideration. It will have a L-3 Communications MUMT-X datalink in place of two older counterparts, communicating on C, D, L, and Ku frequency bands to transmit and receive data and video with all Army UAVs. Lots 5 and 6 will have Link 16 data links. The AH-64E is to be fit for maritime use. The U.S. Army expressed interest in extended-range fuel tanks for greater endurance. , 500 AH-64Es have been delivered. Work on a further upgraded AH-64E, version 6.5 was initiated by the U.S. Army in 2021, and first flew in 2023. The AH-64E is reported to have had problems with its electrical generation systems causing increased scrutiny. AH-64F In 2014, Boeing conceptualized an Apache upgrade prior to the introduction of the U.S. Army's anticipated attack version of the Future Vertical Lift (FVL) aircraft, forecast to replace the Apache by 2040. The conceptual AH-64F would have greater speed via a new 3,000 shp turboshaft engine from the Improved Turbine Engine Program, retractable landing gear, stub wings to offload lift from the main rotor during cruise, and a tail rotor that can articulate 90 degrees to provide forward thrust. In October 2016, the Army revealed they would not pursue another Apache upgrade to focus on funding FVL; the Army will continue buying the Apache through the 2020s until Boeing's production line ends in 2026, then FVL is slated to come online in 2030. Piasecki Speed Apache In the late 1990s, Piasecki conceived a proposal for an Apache with a tail–mounted ducted pusher rotor to enhance speed. However, the Speed Apache was not proceeded with. This appears to be similar to the Piasecki X-49 SpeedHawk. Compound Apache In October 2018, Boeing began testing the AH-64E Block 2 Compound, a compound helicopter design which added a larger fixed wing and a pusher propeller to the Apache airframe to provide additional lift and thrust, respectively. In addition, the engine exhaust was redirected downwards. Collectively, the modifications were anticipated to improve speed to , range to , payload to , and fuel economy. A 30% scale model completed wind tunnel testing in January 2019. The Compound Apache has been pitched as an interim replacement for the Apache before its replacement under the Future Vertical Lift program. Sea Apache During the 1980s naval versions of the AH-64A for the United States Marine Corps and Navy were examined. Multiple concepts were studied with altered landing gear arrangements, improved avionics and weapons. The USMC conducted a two-week evaluation of the Apache in September 1981, including shipboard operation tests. Funding for a naval version was not provided; the USMC continued to use the AH-1. The Canadian Forces Maritime Command also examined naval Apaches. In 2004, British Army AgustaWestland Apaches were deployed upon the Royal Navy's , a Landing Platform Helicopter, for suitability testing; there was U.S. interest in the trials. During the 2011 military intervention in Libya, the British Army extensively used Apaches from HMS Ocean. In 2013, US AH-64Ds were tested on several U.S. Navy ships. Export Apaches Several models have been derived from the AH-64A, AH-64D and AH-64E for export. The British-built AgustaWestland Apache (assembled from kits purchased from Boeing) is based on the AH-64D Block I with several different systems, including more powerful engines, folding rotor blades, and other modifications for operation from Royal Navy vessels. Block modification While a major change in design or role will cause the type designator suffix to change, for example from AH-64D to AH-64E, the helicopters are also subject to block modification. Block modification is the combining of equipment changes into blocks of modification work orders, the modifications in the block (sometimes called a block package) are all done to the helicopter at the same time. Operators The AH-64 has had some modest success on the export market. It has been popular in the middle east, where its role the Gulf War established a bit of reputation. Its main user continues to be the U.S. Army, and the largest user in Europe will be Poland. Australian Army – 29 AH-64Es on order Egyptian Air Force – 46 AH-64Ds Hellenic Army – 28 AH-64A/Ds Indian Air Force – 22 AH-64Es in inventory as of July 2020 Indian Army – 6 AH-64Es on order Indonesian Army 8 AH-64Es Israeli Air Force – 48 AH-64A/Ds Japan Ground Self-Defense Force – 12 AH-64Ds In December 2022, the Japanese government decided to replace its AH-64Ds and other aircraft with unmanned aerial vehicles. Kuwait Air Force – 24 AH-64Ds Royal Moroccan Air Force – 36 AH-64Es on order Royal Netherlands Air Force – 28 AH-64Es Polish Land Forces – 96 AH-64Es on order, 8 AH-64Ds leased Qatar Emiri Air Force – 24 AH-64Es Royal Saudi Land Forces – 47 AH-64A/D/Es Saudi Arabia National Guard – 12 AH-64Es Republic of Singapore Air Force – 19 AH-64Ds Republic of Korea Army – 36 AH-64Es An additional 36 AH-64Es were approved in a $3.5 billion deal. Republic of China Army – 29 AH-64Es United Arab Emirates Air Force – 28 AH-64D/Es British Army – 41 AH-64Es, to be increased to 50 in 2025. United States Army – 819 AH-64D/Es Specifications (AH-64A/D) Notable appearances in media	Technology	Specific aircraft	null
37751	https://en.wikipedia.org/wiki/Turtle	Turtle	Turtles are reptiles of the order Testudines, characterized by a special shell developed mainly from their ribs. Modern turtles are divided into two major groups, the Pleurodira (side necked turtles) and Cryptodira (hidden necked turtles), which differ in the way the head retracts. There are 360 living and recently extinct species of turtles, including land-dwelling tortoises and freshwater terrapins. They are found on most continents, some islands and, in the case of sea turtles, much of the ocean. Like other amniotes (reptiles, birds, and mammals) they breathe air and do not lay eggs underwater, although many species live in or around water. Turtle shells are made mostly of bone; the upper part is the domed carapace, while the underside is the flatter plastron or belly-plate. Its outer surface is covered in scales made of keratin, the material of hair, horns, and claws. The carapace bones develop from ribs that grow sideways and develop into broad flat plates that join up to cover the body. Turtles are ectotherms or "cold-blooded", meaning that their internal temperature varies with their direct environment. They are generally opportunistic omnivores and mainly feed on plants and animals with limited movements. Many turtles migrate short distances seasonally. Sea turtles are the only reptiles that migrate long distances to lay their eggs on a favored beach. Turtles have appeared in myths and folktales around the world. Some terrestrial and freshwater species are widely kept as pets. Turtles have been hunted for their meat, for use in traditional medicine, and for their shells. Sea turtles are often killed accidentally as bycatch in fishing nets. Turtle habitats around the world are being destroyed. As a result of these pressures, many species are extinct or threatened with extinction. Naming and etymology The word turtle is borrowed from the French word or 'turtle, tortoise'. It is a common name and may be used without knowledge of taxonomic distinctions. In North America, it may denote the order as a whole. In Britain, the name is used for sea turtles as opposed to freshwater terrapins and land-dwelling tortoises. In Australia, which lacks true tortoises (family Testudinidae), non-marine turtles were traditionally called tortoises, but more recently turtle has been used for the entire group. The name of the order, Testudines ( ), is based on the Latin word 'tortoise'; and was coined by German naturalist August Batsch in 1788. The order has also been historically known as Chelonii (Latreille 1800) and Chelonia (Ross and Macartney 1802), which are based on the Ancient Greek word () 'tortoise'. Testudines is the official order name due to the principle of priority. The term chelonian is used as a formal name for members of the group. Anatomy and physiology Size The largest living species of turtle (and fourth-largest reptile) is the leatherback turtle, which can reach over in length and weigh over . The largest known turtle was Archelon ischyros, a Late Cretaceous sea turtle up to long, wide between the tips of the front flippers, and estimated to have weighed over . The smallest living turtle is Chersobius signatus of South Africa, measuring no more than in length and weighing . Shell The shell of a turtle is unique among vertebrates and serves to protect the animal and provide shelter from the elements. It is primarily made of 50–60 bones and consists of two parts: the domed, dorsal (back) carapace and the flatter, ventral (belly) plastron. They are connected by lateral (side) extensions of the plastron. The carapace is fused with the vertebrae and ribs while the plastron is formed from bones of the shoulder girdle, sternum, and gastralia (abdominal ribs). During development, the ribs grow sideways into a carapacial ridge, unique to turtles, entering the dermis (inner skin) of the back to support the carapace. The development is signaled locally by proteins known as fibroblast growth factors that include FGF10. The shoulder girdle in turtles is made up of two bones, the scapula and the coracoid. Both the shoulder and pelvic girdles of turtles are located within the shell and hence are effectively within the rib cage. The trunk ribs grow over the shoulder girdle during development. The shell is covered in epidermal (outer skin) scales known as scutes that are made of keratin, the same substance that makes up hair and fingernails. Typically, a turtle has 38 scutes on the carapace and 16 on the plastron, giving them 54 in total. Carapace scutes are divided into "marginals" around the margin and "vertebrals" over the vertebral column, though the scute that overlays the neck is called the "cervical". "Pleurals" are present between the marginals and vertebrals. Plastron scutes include gulars (throat), humerals, pectorals, abdominals, and anals. Side-necked turtles additionally have "intergular" scutes between the gulars. Turtle scutes are usually structured like mosaic tiles, but some species, like the hawksbill sea turtle, have overlapping scutes on the carapace. The shapes of turtle shells vary with the adaptations of the individual species, and sometimes with sex. Land-dwelling turtles are more dome-shaped, which appears to make them more resistant to being crushed by large animals. Aquatic turtles have flatter, smoother shells that allow them to cut through the water. Sea turtles in particular have streamlined shells that reduce drag and increase stability in the open ocean. Some turtle species have pointy or spiked shells that provide extra protection from predators and camouflage against the leafy ground. The lumps of a tortoise shell can tilt its body when it gets flipped over, allowing it to flip back. In male tortoises, the tip of the plastron is thickened and used for butting and ramming during combat. Shells vary in flexibility. Some species, such as box turtles, lack the lateral extensions and instead have the carapace bones fully fused or ankylosed together. Several species have hinges on their shells, usually on the plastron, which allow them to expand and contract. Softshell turtles have rubbery edges, due to the loss of bones. The leatherback turtle has hardly any bones in its shell, but has thick connective tissue and an outer layer of leathery skin. Head and neck The turtle's skull is unique among living amniotes (which includes reptiles, birds and mammals); it is solid and rigid with no openings for muscle attachment (temporal fenestrae). Muscles instead attach to recesses in the back of the skull. Turtle skulls vary in shape, from the long and narrow skulls of softshells to the broad and flattened skull of the mata mata. Some turtle species have developed large and thick heads, allowing for greater muscle mass and stronger bites. Turtles that are carnivorous or durophagous (eating hard-shelled animals) have the most powerful bites. For example, the durophagous Mesoclemmys nasuta has a bite force of . Species that are insectivorous, piscivorous (fish-eating), or omnivorous have lower bite forces. Living turtles lack teeth but have beaks made of keratin sheaths along the edges of the jaws. These sheaths may have sharp edges for cutting meat, serrations for clipping plants, or broad plates for breaking mollusks. Sea turtles, and several extinct forms, have evolved a bony secondary palate which completely separates the oral and nasal cavities. The necks of turtles are highly flexible, possibly to compensate for their rigid shells. Some species, like sea turtles, have short necks while others, such as snake-necked turtles, have long ones. Despite this, all turtle species have eight neck vertebrae, a consistency not found in other reptiles but similar to mammals. Some snake-necked turtles have both long necks and large heads, limiting their ability to lift them when not in water. Some turtles have folded structures in the larynx or glottis that vibrate to produce sound. Other species have elastin-rich vocal cords. Limbs and locomotion Due to their heavy shells, turtles are slow-moving on land. A desert tortoise moves at only . By contrast, sea turtles can swim at . The limbs of turtles are adapted for various means of locomotion and habits and most have five toes. Tortoises are specialized for terrestrial environments and have column-like legs with elephant-like feet and short toes. The gopher tortoise has flattened front limbs for digging in the substrate. Freshwater turtles have more flexible legs and longer toes with webbing, giving them thrust in the water. Some of these species, such as snapping turtles and mud turtles, mainly walk along the water bottom, as they would on land. Others, such as terrapins, swim by paddling with all four limbs, switching between the opposing front and hind limbs, which keeps their direction stable. Sea turtles and the pig-nosed turtle are the most specialized for swimming. Their front limbs have evolved into flippers while the shorter hind limbs are shaped more like rudders. The front limbs provide most of the thrust for swimming, while the hind limbs serve as stabilizers. Sea turtles such as the green sea turtle rotate the front limb flippers like a bird's wings to generate a propulsive force on both the upstroke and on the downstroke. This is in contrast to similar-sized freshwater turtles (measurements having been made on young animals in each case) such as the Caspian turtle, which uses the front limbs like the oars of a rowing boat, creating substantial negative thrust on the recovery stroke in each cycle. In addition, the streamlining of the marine turtles reduces drag. As a result, marine turtles produce a propulsive force twice as large, and swim six times as fast, as freshwater turtles. The swimming efficiency of young marine turtles is similar to that of fast-swimming fish of open water, like mackerel. Compared to other reptiles, turtles tend to have reduced tails, but these vary in both length and thickness among species and between sexes. Snapping turtles and the big-headed turtle have longer tails; the latter uses it for balance while climbing. The cloaca is found underneath and at the base, and the tail itself houses the reproductive organs. Hence, males have longer tails to contain the penis. In sea turtles, the tail is longer and more prehensile in males, who use it to grasp mates. Several turtle species have spines on their tails. Senses Turtles make use of vision to find food and mates, avoid predators, and orient themselves. The retina's light-sensitive cells include both rods for vision in low light, and cones with three different photopigments for bright light, where they have full-color vision. There is possibly a fourth type of cone that detects ultraviolet, as hatchling sea turtles respond experimentally to ultraviolet light, but it is unknown if they can distinguish this from longer wavelengths. A freshwater turtle, the red-eared slider, has an exceptional seven types of cone cell. Sea turtles orient themselves on land by night, using visual features detected in dim light. They can use their eyes in clear surface water, muddy coasts, the darkness of the deep ocean, and also above water. Unlike in terrestrial turtles, the cornea (the curved surface that lets light into the eye) does not help to focus light on the retina, so focusing underwater is handled entirely by the lens, behind the cornea. The cone cells contain oil droplets placed to shift perception toward the red part of the spectrum, improving color discrimination. Visual acuity, studied in hatchlings, is highest in a horizontal band with retinal cells packed about twice as densely as elsewhere. This gives the best vision along the visual horizon. Sea turtles do not appear to use polarized light for orientation as many other animals do. The deep-diving leatherback turtle lacks specific adaptations to low light, such as large eyes, large lenses, or a reflective tapetum. It may rely on seeing the bioluminescence of prey when hunting in deep water. Turtles have no ear openings; the eardrum is covered with scales and encircled by a bony otic capsule, which is absent in other reptiles. Their hearing thresholds are high in comparison to other reptiles, reaching up to 500 Hz in air, but underwater they are more attuned to lower frequencies. The loggerhead sea turtle has been shown experimentally to respond to low sounds, with maximal sensitivity between 100 and 400 Hz. Turtles have olfactory (smell) and vomeronasal receptors along the nasal cavity, the latter of which are used to detect chemical signals. Experiments on green sea turtles showed they could learn to respond to a selection of different odorant chemicals such as triethylamine and cinnamaldehyde, which were detected by olfaction in the nose. Such signals could be used in navigation. Breathing The rigid shell of turtles is not capable of expanding and making room for the lungs, as in other amniotes, so they have had to evolve special adaptations for respiration. The lungs of turtles are attached directly to the carapace above while below, connective tissue attaches them to the organs. They have multiple lateral (side) and medial (middle) chambers (the numbers of which vary between species) and one terminal (end) chamber. The lungs are ventilated using specific groups of abdominal muscles attached to the organs that pull and push on them. Specifically, it is the turtle's large liver that compresses the lungs. Underneath the lungs, in the coelomic cavity, the liver is connected to the right lung by the root, and the stomach is directly attached to the left lung, and to the liver by a mesentery. When the liver is pulled down, inhalation begins. Supporting the lungs is a wall or septum, which is thought to prevent them from collapsing. During exhalation, the contraction of the transversus abdominis muscle propels the organs into the lungs and expels air. Conversely, during inhalation, the relaxing and flattening of the oblique abdominis muscle pulls the transversus back down, allowing air back into the lungs. Although many turtles spend large amounts of their lives underwater, all turtles breathe air and must surface at regular intervals to refill their lungs. Depending on the species, immersion periods vary between a minute and an hour. Some species can respire through the cloaca, which contains large sacs that are lined with many finger-like projections that take up dissolved oxygen from the water. Circulation Turtles share the linked circulatory and pulmonary (lung) systems of vertebrates, where the three-chambered heart pumps deoxygenated blood through the lungs and then pumps the returned oxygenated blood through the body's tissues. The cardiopulmonary system has both structural and physiological adaptations that distinguish it from other vertebrates. Turtles have a large lung volume and can move blood through non-pulmonary blood vessels, including some within the heart, to avoid the lungs while they are not breathing. They can hold their breath for much longer periods than other reptiles and they can tolerate the resulting low oxygen levels. They can moderate the increase in acidity during anaerobic (non-oxygen-based) respiration by chemical buffering and they can lie dormant for months, in aestivation or brumation. The heart has two atria but only one ventricle. The ventricle is subdivided into three chambers. A muscular ridge enables a complex pattern of blood flow so that the blood can be directed either to the lungs via the pulmonary artery, or to the body via the aorta. The ability to separate the two outflows varies between species. The leatherback has a powerful muscular ridge enabling almost complete separation of the outflows, supporting its actively swimming lifestyle. The ridge is less well developed in freshwater turtles like the sliders (Trachemys). Turtles are capable of enduring periods of anaerobic respiration longer than many other vertebrates. This process breaks down sugars incompletely to lactic acid, rather than all the way to carbon dioxide and water as in aerobic (oxygen-based) respiration. They make use of the shell as a source of additional buffering agents for combating increased acidity, and as a sink for lactic acid. Osmoregulation In sea turtles, the bladder is one unit and in most freshwater turtles, it is double-lobed. Sea turtle bladders are connected to two small accessory bladders, located at the sides to the neck of the urinary bladder and above the pubis. Arid-living tortoises have bladders that serve as reserves of water, storing up to 20% of their body weight in fluids. The fluids are normally low in solutes, but higher during droughts when the reptile gains potassium salts from its plant diet. The bladder stores these salts until the tortoise finds fresh drinking water. To regulate the amount of salt in their bodies, sea turtles and the brackish-living diamondback terrapin secrete excess salt in a thick sticky substance from their tear glands. Because of this, sea turtles may appear to be "crying" when on land. Thermoregulation Turtles, like other reptiles, have a limited ability to regulate their body temperature. This ability varies between species, and with body size. Small pond turtles regulate their temperature by crawling out of the water and basking in the sun, while small terrestrial turtles move between sunny and shady places to adjust their temperature. Large species, both terrestrial and marine, have sufficient mass to give them substantial thermal inertia, meaning that they heat up or cool down over many hours. The Aldabra giant tortoise weighs up to some and is able to allow its temperature to rise to some on a hot day, and to fall naturally to around by night. Some giant tortoises seek out shade to avoid overheating on sunny days. On Grand Terre Island, food is scarce inland, shade is scarce near the coast, and the tortoises compete for space under the few trees on hot days. Large males may push smaller females out of the shade, and some then overheat and die. Adult sea turtles, too, have large enough bodies that they can to some extent control their temperature. The largest turtle, the leatherback, can swim in the waters off Nova Scotia, which may be as cold as , while their body temperature has been measured at up to warmer than the surrounding water. To help keep their temperature up, they have a system of countercurrent heat exchange in the blood vessels between their body core and the skin of their flippers. The vessels supplying the head are insulated by fat around the neck. Behavior Diet and feeding Most turtle species are opportunistic omnivores; land-dwelling species are more herbivorous and aquatic ones more carnivorous. Generally lacking speed and agility, most turtles feed either on plant material or on animals with limited movements like mollusks, worms, and insect larvae. Some species, such as the African helmeted turtle and snapping turtles, eat fish, amphibians, reptiles (including other turtles), birds, and mammals. They may take them by ambush but also scavenge. The alligator snapping turtle has a worm-like appendage on its tongue that it uses to lure fish into its mouth. Tortoises are the most herbivorous group, consuming grasses, leaves, and fruits. Many turtle species, including tortoises, supplement their diet with eggshells, animal bones, hair, and droppings for extra nutrients. Turtles generally eat their food in a straightforward way, though some species have special feeding techniques. The yellow-spotted river turtle and the painted turtle may filter feed by skimming the water surface with their mouth and throat open to collect particles of food. When the mouth closes, the throat constricts and water is pushed out through the nostrils and the gap in between the jaws. Some species employ a "gape-and-suck method" where the turtle opens its jaws and expands its throat widely, sucking the prey in. The diet of an individual within a species may change with age, sex, and season, and may also differ between populations. In many species, juveniles are generally carnivorous but become more herbivorous as adults. With Barbour's map turtle, the larger female mainly eats mollusks while the male usually eats arthropods. Blanding's turtle may feed mainly on snails or crayfish depending on the population. The European pond turtle has been recorded as being mostly carnivorous much of the year but switching to water lilies during the summer. Some species have developed specialized diets such as the hawksbill, which eats sponges, the leatherback, which feeds on jellyfish, and the Mekong snail-eating turtle. Communication and intelligence While popularly thought of as mute, turtles make various sounds to communicate. One study which recorded 53 species found that all of them vocalized. Tortoises may bellow when courting and mating. Various species of both freshwater and sea turtles emit short, low-frequency calls from the time they are in the egg to when they are adults. These vocalizations may serve to create group cohesion when migrating. The oblong turtle has a particularly large vocal range; producing sounds described as clacks, clicks, squawks, hoots, various kinds of chirps, wails, , grunts, growls, blow bursts, howls, and drum rolls. Play behavior has been documented in some turtle species. In the laboratory, Florida red-bellied cooters can learn novel tasks and have demonstrated a long-term memory of at least 7.5 months. Similarly, giant tortoises can learn and remember tasks, and master lessons much faster when trained in groups. Tortoises appear to be able to retain operant conditioning nine years after their initial training. Studies have shown that turtles can navigate the environment using landmarks and a map-like system resulting in accurate direct routes towards a goal. Navigation in turtles have been correlated to high cognition function in the medial cortex region of the brain. Defense When sensing danger, a turtle may flee, freeze or withdraw into its shell. Freshwater turtles flee into the water, though the Sonora mud turtle may take refuge on land as the shallow temporary ponds they inhabit make them vulnerable. When startled, a softshell turtle may dive underwater and bury itself under the sea floor. If a predator persists, the turtle may bite or discharge from its cloaca. Several species produce foul-smelling chemicals from musk glands. Other tactics include threat displays and Bell's hinge-back tortoise can play dead. When attacked, big-headed turtle hatchlings squeal, possibly startling the predator. Migration Turtles are the only reptiles that migrate long distances, more specifically the marine species that can travel up to thousands of kilometers. Some non-marine turtles, such as the species of Geochelone (terrestrial), Chelydra (freshwater), and Malaclemys (estuarine), migrate seasonally over much shorter distances, up to around , to lay eggs. Such short migrations are comparable to those of some lizards, snakes, and crocodilians. Sea turtles nest in a specific area, such as a beach, leaving the eggs to hatch unattended. The young turtles leave that area, migrating long distances in the years or decades in which they grow to maturity, and then return seemingly to the same area every few years to mate and lay eggs, though the precision varies between species and populations. This "natal homing" has appeared remarkable to biologists, though there is now plentiful evidence for it, including from genetics. How sea turtles navigate to their breeding beaches remains unknown. One possibility is imprinting as in salmon, where the young learn the chemical signature, effectively the scent, of their home waters before leaving, and remember that when the time comes for them to return as adults. Another possible cue is the orientation of the Earth's magnetic field at the natal beach. There is experimental evidence that turtles have an effective magnetic sense, and that they use this in navigation. Proof that homing occurs is derived from genetic analysis of populations of loggerheads, hawksbills, leatherbacks, and olive ridleys by nesting place. For each of these species, the populations in different places have their own mitochondrial DNA genetic signatures that persist over the years. This shows that the populations are distinct and that homing must be occurring reliably. Reproduction and life cycle Turtles have a wide variety of mating behaviors but do not form pair-bonds or social groups. In green sea turtles, females generally outnumber males. In terrestrial species, males are often larger than females and fighting between males establishes a dominance hierarchy for access to mates. For most semi-aquatic and bottom-walking aquatic species, combat occurs less often. Males of these species instead may use their size advantage to mate forcibly. In fully aquatic species, males are often smaller than females and rely on courtship displays to gain mating access to females. Courtship and mounting Courtship varies between species, and with habitat. It is often complex in aquatic species, both marine and freshwater, but simpler in the semi-aquatic mud turtles and snapping turtles. A male tortoise bobs his head, then subdues the female by biting and butting her before mounting. The male scorpion mud turtle approaches the female from the rear, and often resorts to aggressive methods such as biting the female's tail or hind limbs, followed by a mounting. Female choice is important in some species, and female green sea turtles are not always receptive. As such, they have evolved behaviors to avoid the male's attempts at copulation, such as swimming away, confronting the male followed by biting or taking up a refusal position with her body vertical, her limbs widely outspread, and her plastron facing the male. If the water is too shallow for the refusal position, the females resort to beaching themselves, as the males do not follow them ashore. All turtles fertilize internally; mounting and copulation can be difficult. In many species, males have a concave plastron that interlocks with the female's carapace. In species like the Russian tortoise, the male has a lighter shell and longer legs. The high, rounded shape of box turtles are particular obstacles for mounting. The male eastern box turtle leans backward and hooks onto the back of the female's plastron. Aquatic turtles mount in water, and female sea turtles support the mounting male while swimming and diving. During copulation, the male turtle aligns his tail with the female's so he can insert his penis into her cloaca. Some female turtles can store sperm from multiple males and their egg clutches can have multiple sires. Eggs and hatchlings Turtles, including sea turtles, lay their eggs on land, although some lay eggs near water that rises and falls in level, submerging the eggs. While most species build nests and lay eggs where they forage, some travel miles. The common snapping turtle walks on land, while sea turtles travel even further; the leatherback swims some to its nesting beaches. Most turtles create a nest for their eggs. Females usually dig a flask-like chamber in the substrate. Other species lay their eggs in vegetation or crevices. Females choose nesting locations based on environmental factors such as temperature and humidity, which are important for developing embryos. Depending on the species, the number of eggs laid varies from one to over 100. Larger females can lay eggs that are greater in number or bigger in size. Compared to freshwater turtles, tortoises deposit fewer but larger eggs. Females can lay multiple clutches throughout a season, particularly in species that experience unpredictable monsoons. Most mother turtles do no more in the way of parental care than covering their eggs and immediately leaving, though some species guard their nests for days or weeks. Eggs vary between rounded, oval, elongated, and between hard- and soft-shelled. Most species have their sex determined by temperature. In some species, higher temperatures produce females and lower ones produce males, while in others, milder temperatures produce males and both hot and cold extremes produce females. There is experimental evidence that the embryos of Mauremys reevesii can move around inside their eggs to select the best temperature for development, thus influencing their sexual destiny. In other species, sex is determined genetically. The length of incubation for turtle eggs varies from two to three months for temperate species, and four months to over a year for tropical species. Species that live in warm temperate climates can delay their development. Hatching young turtles break out of the shell using an egg tooth, a sharp projection that exists temporarily on their upper beak. Hatchlings dig themselves out of the nest and find safety in vegetation or water. Some species stay in the nest for longer, be it for overwintering or to wait for the rain to loosen the soil for them to dig out. Young turtles are highly vulnerable to predators, both in the egg and as hatchlings. Mortality is high during this period but significantly decreases when they reach adulthood. Most species grow quickly during their early years and slow down when they mature. Lifespan Turtles can live long lives. The oldest living turtle and land animal is said to be a Seychelles giant tortoise named Jonathan, who turned 187 in 2019. A Galápagos tortoise named Harriet was collected by Charles Darwin in 1835; it died in 2006, having lived for at least 176 years. Most wild turtles do not reach that age. Turtles keep growing new scutes under the previous scutes every year, allowing researchers to estimate how long they have lived. They also age slowly. The survival rate for adult turtles can reach 99% per year. Systematics and evolution Fossil history Zoologists have sought to explain the evolutionary origin of the turtles, and in particular of their unique shells. In 1914, Jan Versluys proposed that bony plates in the dermis, called osteoderms, fused to the ribs beneath them, later called the "Polka Dot Ancestor" by Olivier Rieppel. The theory accounted for the evolution of fossil pareiasaurs from Bradysaurus to Anthodon, but not for how the ribs could have become attached to the bony dermal plates. More recent discoveries have painted a different scenario for the evolution of the turtle's shell. The stem-turtles Eunotosaurus of the Middle Permian, Pappochelys of the Middle Triassic, and Eorhynchochelys of the Late Triassic lacked carapaces and plastrons but had shortened torsos, expanded ribs, and lengthened dorsal vertebrae. Also in the Late Triassic, Odontochelys had a partial shell consisting of a complete bony plastron and an incomplete carapace. The development of a shell reached completion with the Late Triassic Proganochelys, with its fully developed carapace and plastron. Adaptations that led to the evolution of the shell may have originally been for digging and a fossorial lifestyle. The oldest known members of the Pleurodira lineage are the Platychelyidae, from the Late Jurassic. The oldest known unambiguous cryptodire is Sinaspideretes, a close relative of softshell turtles, from the Late Jurassic of China. Turtles became highly diverse during the Cretaceous, as climatic conditions in this period were favourable for their global dispersal. During the Late Cretaceous and Cenozoic, members of the pleurodire families Bothremydidae and Podocnemididae became widely distributed in the Northern Hemisphere due to their coastal habits. The oldest known soft-shelled turtles and sea turtles appeared during the Early Cretaceous. Tortoises originated in Asia during the Eocene. A late surviving group of stem-turtles, the Meiolaniidae, survived in Australasia into the Pleistocene and Holocene. External relationships The turtles' exact ancestry has been disputed. It was believed they were the only surviving branch of the ancient evolutionary grade Anapsida, which includes groups such as procolophonids and pareiasaurs. All anapsid skulls lack a temporal opening while all other living amniotes have temporal openings. It was later suggested that the anapsid-like turtle skulls may be due to backward evolution rather than to anapsid descent. Fossil evidence has shown that early stem-turtles possessed small temporal openings. Some early morphological phylogenetic studies have placed turtles closer to Lepidosauria (tuataras, lizards, and snakes) than to Archosauria (crocodilians and birds). By contrast, several molecular studies place turtles either within Archosauria, or, more commonly, as a sister group to extant archosaurs, though an analysis conducted by Tyler Lyson and colleagues (2012) recovered turtles as the sister group of lepidosaurs instead. Ylenia Chiari and colleagues (2012) analyzed 248 nuclear genes from 16 vertebrates and suggested that turtles share a more recent common ancestor with birds and crocodilians. The date of separation of turtles and birds and crocodilians was estimated to be during the Permian. Through genomic-scale phylogenetic study of ultra-conserved elements (UCEs) to clarify the placement of turtles within reptiles, Nicholas Crawford and colleagues (2012) similarly found that turtles are closer to birds and crocodilians. Using the draft (unfinished) genome sequences of the green sea turtle and the Chinese softshell turtle, Zhuo Wang and colleagues (2013) concluded that turtles are likely a sister group of crocodilians and birds. The external phylogeny of the turtles is shown in the cladogram below. Internal relationships Modern turtles and their extinct relatives with a complete shell are classified within the clade Testudinata. The most recent common ancestor of living turtles, corresponding to the split between Pleurodira (side-necked species) and Cryptodira (hidden necked species), is estimated to have occurred around during the Late Triassic. Robert Thompson and colleagues (2021) comment that living turtles have low diversity, relative to how long they existed. Diversity has been stable, according to their analysis, except for a single rapid increase around the Eocene-Oligocene boundary some 30 million years ago, and a large regional extinction at roughly the same time. They suggest that global climate change caused both events, as the cooling and drying caused the land to become arid and turtles to become extinct there, while new continental margins opened up by the climate change provided habitats for other species to evolve. The cladogram, from Nicholas Crawford and colleagues 2015, shows the internal phylogeny of the Testudines down to the level of families. The analysis by Thompson and colleagues in 2021 supports the same structure down to the family level. Differences between the two suborders Turtles are divided into two living suborders: Cryptodira and Pleurodira. The two groups differ in the way the neck is retracted for protection. Pleurodirans retract their neck to the side and in front of the shoulder girdles, whereas cryptodirans retract their neck backward into their shell. These motions are enabled by the morphology and arrangement of neck vertebrae. Sea turtles (which belong to Cryptodira) have mostly lost the ability to retract their heads. The adductor muscles in the lower jaw create a pulley-like system in both subgroups. However, the bones that the muscles articulate with differ. In Pleurodira, the pulley is formed with the pterygoid bones of the palate, but in Cryptodira the pulley is formed with the otic capsule. Both systems help to vertically redirect the adductor muscles and maintain a powerful bite. A further difference between the suborders is the attachment of the pelvis. In Cryptodira, the pelvis is free, linked to the shell only by ligaments. In Pleurodira, the pelvis is sutured, joined with bony connections, to the carapace and to the plastron, creating a pair of large columns of bone at the back end of the turtle, linking the two parts of the shell. Distribution and habitat Turtles are widely distributed across the world's continents, oceans, and islands with terrestrial, fully aquatic, and semi-aquatic species. Sea turtles are mainly tropical and subtropical, but leatherbacks can be found in colder areas of the Atlantic and Pacific. Living Pleurodira all live in freshwater and are found only in the Southern Hemisphere. The Cryptodira include terrestrial, freshwater, and marine species, and these range more widely. The world regions richest in non-marine turtle species are the Amazon basin, the Gulf of Mexico drainages of the United States, and parts of South and Southeast Asia. For turtles in colder climates, their distribution is limited by constraints on reproduction, which is reduced by long hibernations. North American species barely range above the southern Canadian border. Some turtles are found at high altitudes, for example, the species Terrapene ornata occurs up to in New Mexico. Conversely, the leatherback sea turtle can dive over . Species of the genus Gopherus can tolerate both below freezing and over in body temperature, though they are most active at . Conservation Among vertebrate orders, turtles are second only to primates in the percentage of threatened species. 360 modern species have existed since 1500 AD. Of these, 51–56% are considered threatened and 60% considered threatened or extinct. Turtles face many threats, including habitat destruction, harvesting for consumption, the pet trade, turtle racing, light pollution, and climate change. Asian species have a particularly high extinction risk, primarily due to their long-term unsustainable exploitation for food and medicine, and about 83% of Asia's non-marine turtle species are considered threatened. As of 2021, turtle extinction is progressing much faster than during the Cretaceous-Tertiary extinction. At this rate, all turtles could be extinct in a few centuries. Turtle hatcheries can be set up when protection against flooding, erosion, predation, or heavy poaching is required. Chinese markets have sought to satisfy an increasing demand for turtle meat with farmed turtles. In 2007 it was estimated that over a thousand turtle farms operated in China. All the same, wild turtles continue to be caught and sent to market in large numbers, resulting in what conservationists have called "the Asian turtle crisis". In the words of the biologist George Amato, the hunting of turtles "vacuumed up entire species from areas in Southeast Asia", even as biologists still did not know how many species lived in the region. In 2000, all the Asian box turtles were placed on the CITES list of endangered species. Harvesting wild turtles is legal in some American states, and there has been a growing demand for American turtles in China. The Florida Fish and Wildlife Conservation Commission estimated in 2008 that around 3,000 pounds of softshell turtles were exported weekly via Tampa International Airport. However, the great majority of turtles exported from the US between 2002 and 2005 were farmed. Large numbers of sea turtles are accidentally killed in longlines, gillnets, and trawling nets as bycatch. A 2010 study suggested that over 8 million had been killed between 1990 and 2008; the Eastern Pacific and the Mediterranean were identified as among the areas worst affected. Since the 1980s, the United States has required all shrimp trawlers to fit their nets with turtle excluder devices that prevent turtles from being entangled in the net and drowning. More locally, other human activities are affecting marine turtles. In Australia, Queensland's shark culling program, which uses shark nets and drum lines, has killed over 5,000 turtles as bycatch between 1962 and 2015; including 719 loggerhead turtles and 33 hawksbill sea turtles, which are listed as critically endangered. Native turtle populations can also be threatened by invasive ones. The central North American red-eared slider turtle has been listed among the "world's worst invasive species", pet turtle having been released globally. They appear to compete with native turtle species in eastern and western North America, Europe, and Japan. Human uses On space flights Two tortoises were on the Soviet Union's September 1968 Zond 5 circumlunar flight, making them the first earthly living things to travel to the vicinity of the Moon. Turtles were also on the Zond 6 (1968) and the Zond 7 (1969) circumlunar flights. In culture Turtles have featured in human cultures across the world since ancient times. They are generally viewed positively despite not being "cuddly" or flashy; their association with the ancient times and old age have contributed to their endearing image. In Hindu mythology, the World Turtle, named Kurma or Kacchapa, supports four elephants on his back; they, in turn, carry the weight of the whole world on their backs. The turtle is one of the ten avatars or incarnations of the god Vishnu. The yoga pose Kurmasana is named for the avatar. World Turtles are found in Native American cultures including the Algonquian, Iroquois, and Lenape. They tell many versions of the creation story of Turtle Island. One version has Muskrat pile up earth on Turtle's back, creating the continent of North America. An Iroquois version has the pregnant Sky Woman fall through a hole in the sky between a tree's roots, where she is caught by birds who land her safely on Turtle's back; the Earth grows around her. The turtle here is altruistic, but the world is a heavy burden, and the turtle sometimes shakes itself to relieve the load, causing earthquakes. A turtle was the symbol of the Ancient Mesopotamian god Enki from the 3rd millennium BCE onward. An ancient Greek origin myth told that only the tortoise refused the invitation of the gods Zeus and Hera to their wedding, as it preferred to stay at home. Zeus then ordered it to carry its house with it, ever after. Another of their gods, Hermes, invented a seven-stringed lyre made with the shell of a tortoise. In the Shang dynasty Chinese practice of plastromancy, dating back to 1200 BCE, oracles were obtained by inscribing questions on turtle plastrons using the oldest known form of Chinese characters, burning the plastron, and interpreting the resulting cracks. Later, the turtle was one of the four sacred animals in Confucianism, while in the Han period, steles were mounted on top of stone turtles, later linked with Bixi, the turtle-shelled son of the Dragon King. Marine turtles feature significantly in Australian Aboriginal art. The army of Ancient Rome used the ("tortoise") formation where soldiers would form a shield wall for protection. In Aesop's Fables, "The Tortoise and the Hare" tells how an unequal race may be won by the slower partner. Lewis Carroll's 1865 Alice's Adventures in Wonderland features a Mock Turtle, named for a soup meant to imitate the expensive soup made from real turtle meat. In 1896, the French playwright Léon Gandillot wrote a comedy in three acts named that was "a Parisian sensation" in its run in France, and came to the Manhattan Theatre, Broadway, New York, in 1898 as The Turtle. A "cosmic turtle" and the island motif reappear in Gary Snyder's 1974 novel Turtle Island, and again in Terry Pratchett's Discworld series as Great A'Tuin, starting with the 1983 novel The Colour of Magic. It is supposedly of the species Chelys galactica, the galactic turtle, complete with four elephants on its back to support Discworld. A giant fire-breathing turtle called Gamera is the star of a series of Japanese monster movies in the kaiju genre and has had twelve films from 1965 to 2006. Turtles have been featured in comic books and animations such as the 1984 Teenage Mutant Ninja Turtles. As pets Some turtles, particularly small terrestrial and freshwater species, are kept as pets. The demand for pet turtles increased in the 1950s, with the US being the main supplier, particularly of farm-bred red-eared sliders. The popularity for exotic pets has led to an increase in illegal wildlife trafficking. Around 21% of the value of live animal trade is in reptiles, and turtles are among the more popularly traded species. Poor husbandry of tortoises can cause chronic rhinitis (nasal swelling), overgrown beaks, hyperparathyroidism (which softens their skeleton), constipation, various reproductive problems, and injuries from dogs. In the early 20th century, people in the United States have organized and gambled on turtle races. As food and other uses The flesh of captured wild turtles continues to be eaten in Asian cultures, while turtle soup was once a popular dish in English cuisine. Gopher tortoise stew has been popular with some groups in Florida. The supposed aphrodisiac or medicinal properties of turtle eggs created a large trade for them in Southeast Asia. Hard-shell turtle plastrons and soft-shell carapaces are widely used in traditional Chinese medicine; Taiwan imported nearly 200 metric tons of hard-shells from its neighbors yearly from 1999 to 2008. A popular medicinal preparation based on herbs and turtle shells is guilinggao jelly. The substance tortoiseshell, usually from the hawksbill turtle, has been used for centuries to make jewelry, tools, and ornaments around the Western Pacific. Hawksbills have accordingly been hunted for their shells. The trading of tortoiseshell was internationally banned in 1977 by CITES. Some cultures have used turtle shells to make music: Native American shamans made them into ceremonial rattles, while Aztecs, Mayas, and Mixtecs made drums.	Biology and health sciences	Reptiles	null
37754	https://en.wikipedia.org/wiki/Mountain	Mountain	A mountain is an elevated portion of the Earth's crust, generally with steep sides that show significant exposed bedrock. Although definitions vary, a mountain may differ from a plateau in having a limited summit area, and is usually higher than a hill, typically rising at least above the surrounding land. A few mountains are isolated summits, but most occur in mountain ranges. Mountains are formed through tectonic forces, erosion, or volcanism, which act on time scales of up to tens of millions of years. Once mountain building ceases, mountains are slowly leveled through the action of weathering, through slumping and other forms of mass wasting, as well as through erosion by rivers and glaciers. High elevations on mountains produce colder climates than at sea level at similar latitude. These colder climates strongly affect the ecosystems of mountains: different elevations have different plants and animals. Because of the less hospitable terrain and climate, mountains tend to be used less for agriculture and more for resource extraction, such as mining and logging, along with recreation, such as mountain climbing and skiing. The highest mountain on Earth is Mount Everest in the Himalayas of Asia, whose summit is above mean sea level. The highest known mountain on any planet in the Solar System is Olympus Mons on Mars at . The tallest mountain including submarine terrain is Mauna Kea in Hawaii from its underwater base at 9,330 m (30,610 ft); some scientists consider it to be the tallest on earth. Definition There is no universally accepted definition of a mountain. Elevation, volume, relief, steepness, spacing and continuity have been used as criteria for defining a mountain. In the Oxford English Dictionary a mountain is defined as "a natural elevation of the earth surface rising more or less abruptly from the surrounding level and attaining an altitude which, relatively to the adjacent elevation, is impressive or notable." Whether a landform is called a mountain may depend on local usage. John Whittow's Dictionary of Physical Geography states "Some authorities regard eminences above as mountains, those below being referred to as hills." In the United Kingdom and the Republic of Ireland, a mountain is usually defined as any summit at least high, which accords with the official UK government's definition that a mountain, for the purposes of access, is a summit of or higher. In addition, some definitions also include a topographical prominence requirement, such as that the mountain rises above the surrounding terrain. At one time, the United States Board on Geographic Names defined a mountain as being or taller, but has abandoned the definition since the 1970s. Any similar landform lower than this height was considered a hill. However, today, the United States Geological Survey concludes that these terms do not have technical definitions in the US. The UN Environmental Programme's definition of "mountainous environment" includes any of the following: Class 1: Elevation greater than . Class 2: Elevation between . Class 3: Elevation between . Class 4: Elevation between , with a slope greater than 2 degrees. Class 5: Elevation between , with a slope greater than 5 degrees or elevation range within . Class 6: Elevation between , with a elevation range within . Class 7: Isolated inner basins and plateaus less than in area that are completely surrounded by Class 1 to 6 mountains, but do not themselves meet criteria for Class 1 to 6 mountains. Using these definitions, mountains cover 33% of Eurasia, 19% of South America, 24% of North America, and 14% of Africa. As a whole, 24% of the Earth's land mass is mountainous. Geology There are three main types of mountains: volcanic, fold, and block. All three types are formed from plate tectonics: when portions of the Earth's crust move, crumple, and dive. Compressional forces, isostatic uplift and intrusion of igneous matter forces surface rock upward, creating a landform higher than the surrounding features. The height of the feature makes it either a hill or, if higher and steeper, a mountain. Major mountains tend to occur in long linear arcs, indicating tectonic plate boundaries and activity. Volcanoes Volcanoes are formed when a plate is pushed below another plate, or at a mid-ocean ridge or hotspot. At a depth of around , melting occurs in rock above the slab (due to the addition of water), and forms magma that reaches the surface. When the magma reaches the surface, it often builds a volcanic mountain, such as a shield volcano or a stratovolcano. Examples of volcanoes include Mount Fuji in Japan and Mount Pinatubo in the Philippines. The magma does not have to reach the surface in order to create a mountain: magma that solidifies below ground can still form dome mountains, such as Navajo Mountain in the US. Fold mountains Fold mountains occur when two plates collide: shortening occurs along thrust faults and the crust is overthickened. Since the less dense continental crust "floats" on the denser mantle rocks beneath, the weight of any crustal material forced upward to form hills, plateaus or mountains must be balanced by the buoyancy force of a much greater volume forced downward into the mantle. Thus the continental crust is normally much thicker under mountains, compared to lower lying areas. Rock can fold either symmetrically or asymmetrically. The upfolds are anticlines and the downfolds are synclines: in asymmetric folding there may also be recumbent and overturned folds. The Balkan Mountains and the Jura Mountains are examples of fold mountains. Block mountains Block mountains are caused by faults in the crust: a plane where rocks have moved past each other. When rocks on one side of a fault rise relative to the other, it can form a mountain. The uplifted blocks are block mountains or horsts. The intervening dropped blocks are termed graben: these can be small or form extensive rift valley systems. This kind of landscape can be seen in East Africa, the Vosges and Rhine valley, and the Basin and Range Province of Western North America. These areas often occur when the regional stress is extensional and the crust is thinned. Erosion During and following uplift, mountains are subjected to the agents of erosion (water, wind, ice, and gravity) which gradually wear the uplifted area down. Erosion causes the surface of mountains to be younger than the rocks that form the mountains themselves. Glacial processes produce characteristic landforms, such as pyramidal peaks, knife-edge arêtes, and bowl-shaped cirques that can contain lakes. Plateau mountains, such as the Catskills, are formed from the erosion of an uplifted plateau. Climate Climate in the mountains becomes colder at high elevations, due to an interaction between radiation and convection. Sunlight in the visible spectrum hits the ground and heats it. The ground then heats the air at the surface. If radiation were the only way to transfer heat from the ground to space, the greenhouse effect of gases in the atmosphere would keep the ground at roughly , and the temperature would decay exponentially with height. However, when air is hot, it tends to expand, which lowers its density. Thus, hot air tends to rise and transfer heat upward. This is the process of convection. Convection comes to equilibrium when a parcel of air at a given altitude has the same density as its surroundings. Air is a poor conductor of heat, so a parcel of air will rise and fall without exchanging heat. This is known as an adiabatic process, which has a characteristic pressure-temperature dependence. As the pressure gets lower, the temperature decreases. The rate of decrease of temperature with elevation is known as the adiabatic lapse rate, which is approximately 9.8 °C per kilometre (or per 1000 feet) of altitude. The presence of water in the atmosphere complicates the process of convection. Water vapor contains latent heat of vaporization. As air rises and cools, it eventually becomes saturated and cannot hold its quantity of water vapor. The water vapor condenses to form clouds and releases heat, which changes the lapse rate from the dry adiabatic lapse rate to the moist adiabatic lapse rate (5.5 °C per kilometre or per 1000 feet) The actual lapse rate can vary by altitude and by location. Therefore, moving up on a mountain is roughly equivalent to moving 80 kilometres (45 miles or 0.75° of latitude) towards the nearest pole. This relationship is only approximate, however, since local factors such as proximity to oceans (such as the Arctic Ocean) can drastically modify the climate. As the altitude increases, the main form of precipitation becomes snow and the winds increase. The effect of the climate on the ecology at an elevation can be largely captured through a combination of amount of precipitation, and the biotemperature, as described by Leslie Holdridge in 1947. Biotemperature is the mean temperature; all temperatures below are considered to be 0 °C. When the temperature is below 0 °C, plants are dormant, so the exact temperature is unimportant. The peaks of mountains with permanent snow can have a biotemperature below . Climate change Mountain environments are particularly sensitive to anthropogenic climate change and are currently undergoing alterations unprecedented in last 10,000 years. The effect of global warming on mountain regions (relative to lowlands) is still an active area of study. Observational studies show that highlands are warming faster than nearby lowlands, but when compared globally, the effect disappears. Precipitation in highland areas is not increasing as quickly as in lowland areas. Climate modeling give mixed signals about whether a particular highland area will have increased or decreased precipitation. Climate change has started to affect the physical and ecological systems of mountains. In recent decades mountain ice caps and glaciers have experienced accelerating ice loss. The melting of the glaciers, permafrost and snow has caused underlying surfaces to become increasingly unstable. Landslip hazards have increased in both number and magnitude due to climate change. Patterns of river discharge will also be significantly affected by climate change, which in turn will have significant impacts on communities that rely on water fed from alpine sources. Nearly half of mountain areas provide essential or supportive water resources for mainly urban populations, in particular during the dry season and in semiarid areas such as in central Asia. Alpine ecosystems can be particularly climatically sensitive. Many mid-latitude mountains act as cold climate refugia, with the ecosystems occupying small environmental niches. As well as the direct influence that the change in climate can have on an ecosystem, there is also the indirect one on the soils from changes in stability and soil development. Ecology The colder climate on mountains affects the plants and animals residing on mountains. A particular set of plants and animals tend to be adapted to a relatively narrow range of climate. Thus, ecosystems tend to lie along elevation bands of roughly constant climate. This is called altitudinal zonation. In regions with dry climates, the tendency of mountains to have higher precipitation as well as lower temperatures also provides for varying conditions, which enhances zonation. Some plants and animals found in altitudinal zones tend to become isolated since the conditions above and below a particular zone will be inhospitable and thus constrain their movements or dispersal. These isolated ecological systems are known as sky islands. Altitudinal zones tend to follow a typical pattern. At the highest elevations, trees cannot grow, and whatever life may be present will be of the alpine type, resembling tundra. Just below the tree line, one may find subalpine forests of needleleaf trees, which can withstand cold, dry conditions. Below that, montane forests grow. In the temperate portions of the earth, those forests tend to be needleleaf trees, while in the tropics, they can be broadleaf trees growing in a rainforest. Mountains and humans The highest known permanently tolerable altitude is at . At very high altitudes, the decreasing atmospheric pressure means that less oxygen is available for breathing, and there is less protection against solar radiation (UV). Above elevation, there is not enough oxygen to support human life. This is sometimes referred to as the "death zone". The summits of Mount Everest and K2 are in the death zone. Mountain societies and economies Mountains are generally less preferable for human habitation than lowlands, because of harsh weather and little level ground suitable for agriculture. While 7% of the land area of Earth is above , only 140 million people live above that altitude and only 20-30 million people above elevation. About half of mountain dwellers live in the Andes, Central Asia, and Africa. With limited access to infrastructure, only a handful of human communities exist above of elevation. Many are small and have heavily specialized economies, often relying on industries such as agriculture, mining, and tourism. An example of such a specialized town is La Rinconada, Peru, a gold-mining town and the highest elevation human habitation at . A counterexample is El Alto, Bolivia, at , which has a highly diverse service and manufacturing economy and a population of nearly 1 million. Traditional mountain societies rely on agriculture, with higher risk of crop failure than at lower elevations. Minerals often occur in mountains, with mining being an important component of the economics of some mountain-based societies. More recently, tourism has become more important to the economies of mountain communities, with developments focused around attractions such as national parks and ski resorts. Approximately 80% of mountain people live below the poverty line. Most of the world's rivers are fed from mountain sources, with snow acting as a storage mechanism for downstream users. More than half of humanity depends on mountains for water. In geopolitics, mountains are often seen as natural boundaries between polities. Mountaineering Mountains as sacred places Mountains often play a significant role in religion. There are for example a number of sacred mountains within Greece such as Mount Olympus which was held to be the home of the gods. In Japanese culture, the volcano of Mount Fuji is also held to be sacred with tens of thousands of Japanese ascending it each year. Mount Kailash, in the Tibet Autonomous Region of China, is considered to be sacred in four religions: Hinduism, Bon, Buddhism, and Jainism. In Ireland, pilgrimages are made up the Mount Brandon by Irish Catholics. The Himalayan peak of Nanda Devi is associated with the Hindu goddesses Nanda and Sunanda; it has been off-limits to climbers since 1983. Mount Ararat is a sacred mountain, as it is believed to be the biblical landing place of Noah's Ark. In Europe and especially in the Alps, summit crosses are often erected on the tops of prominent mountains. Superlatives Heights of mountains are typically measured above sea level. Using this metric, Mount Everest is the highest mountain on Earth, at . There are at least 100 mountains with heights of over above sea level, all of which are located in central and southern Asia. The highest mountains above sea level are generally not the highest above the surrounding terrain. There is no precise definition of surrounding base, but Denali, Mount Kilimanjaro and Nanga Parbat are possible candidates for the tallest mountain on land by this measure. The bases of mountain islands are below sea level, and given this consideration Mauna Kea ( above sea level) is the world's tallest mountain and volcano, rising about from the Pacific Ocean floor. The highest mountains are not generally the most voluminous. Mauna Loa () is the largest mountain on Earth in terms of base area (about ) and volume (about ). Mount Kilimanjaro is the largest non-shield volcano in terms of both base area () and volume (). Mount Logan is the largest non-volcanic mountain in base area (). The highest mountains above sea level are also not those with peaks farthest from the centre of the Earth, because the figure of the Earth is not spherical. Sea level closer to the equator is several miles farther from the centre of the Earth. The summit of Chimborazo, Ecuador's tallest mountain, is usually considered to be the farthest point from the Earth's centre, although the southern summit of Peru's tallest mountain, Huascarán, is another contender. Both have elevations above sea level more than less than that of Everest.	Physical sciences	Terrestrial features	null
37764	https://en.wikipedia.org/wiki/Hippopotamus	Hippopotamus	The hippopotamus (Hippopotamus amphibius) (; : hippopotamuses; often shortened to hippo (: hippos), further qualified as the common hippopotamus, Nile hippopotamus and river hippopotamus, is a large semiaquatic mammal native to sub-Saharan Africa. It is one of only two extant species in the family Hippopotamidae, the other being the pygmy hippopotamus (Choeropsis liberiensis or Hexaprotodon liberiensis). Its name comes from the ancient Greek for "river horse" (). After elephants and rhinoceroses, the hippopotamus is the next largest land mammal. It is also the largest extant land artiodactyl. Despite their physical resemblance to pigs and other terrestrial even-toed ungulates, the closest living relatives of the hippopotamids are cetaceans (whales, dolphins, porpoises, etc.), from which they diverged about 55 million years ago. Hippos are recognisable for their barrel-shaped torsos, wide-opening mouths with large canine tusks, nearly hairless bodies, pillar-like legs, and large size: adults average for bulls (males) and for cows (females). Despite its stocky shape and short legs, it is capable of running over short distances. Hippos inhabit rivers, lakes, and mangrove swamps. Territorial bulls each preside over a stretch of water and a group of five to thirty cows and calves. Mating and birth both occur in the water. During the day, hippos remain cool by staying in water or mud, emerging at dusk to graze on grasses. While hippos rest near each other in the water, grazing is a solitary activity and hippos typically do not display territorial behaviour on land. Hippos are among the most dangerous animals in the world due to their aggressive and unpredictable nature. They are threatened by habitat loss and poaching for their meat and ivory (canine teeth). Etymology The Latin word is derived from the ancient Greek (), from () and () , together meaning . In English, the plural is "hippopotamuses". Taxonomy and origins Classification The modern hippopotamus and the pygmy hippopotamus are the only living members of the family Hippopotamidae. Some taxonomists place hippos and anthracotheres in the superfamily Anthracotheroidea. Hippopotamidae are classified along with other even-toed ungulates in the order Artiodactyla. Five subspecies of hippos have been described based on morphological differences in their skulls as well as differences in geographical range: H. a. amphibius – (the nominate subspecies) ranges from Gambia east to Ethiopia and then south to Mozambique and historically ranged as far north as Egypt; its skull is distinguished by a moderately reduced preorbital region, a bulging dorsal surface, elongated mandibular symphysis and larger chewing teeth. H. a. kiboko – found in Kenya and Somalia; was noted to be smaller and more lightly coloured than other hippos with wider nostrils, somewhat longer snout and more rounded and relatively raised orbits with the space between them being incurved. H. a. capensis – found in Zambia and South Africa; distinguished by wider orbits. H. a. tschadensis – ranges between Chad and Niger; featured a slightly shorter but broader face, and pronounced, forward-facing orbits. H. a. constrictus – ranged from the southern Democratic Republic of Congo to Angola and Namibia; skull characterised by a thicker preorbital region, shorter snout, flatter dorsal surface, reduced mandibular symphysis and smaller chewing teeth. The suggested subspecies above were never widely used or validated by field biologists; the described morphological differences were small enough that they could have resulted from simple variation in nonrepresentative samples. A study examining mitochondrial DNA from skin biopsies taken from 13 sampling locations found "low, but significant, genetic differentiation" among H. a. amphibius, H. a. capensis, and H. a. kiboko. Neither H. a. tschadensis nor H. a. constrictus have been tested. Evolution Until 1909, naturalists classified hippos together with pigs based on molar patterns. Several lines of evidence, first from blood proteins, then from molecular systematics, DNA and the fossil record, show their closest living relatives are cetaceans (whales, dolphins, and porpoises). The common ancestor of hippos and whales branched off from Ruminantia and the rest of the even-toed ungulates; the cetacean and hippo lineages split soon afterwards. The most recent theory of the origins of Hippopotamidae suggests hippos and whales shared a common semiaquatic ancestor that branched off from other artiodactyls around . This hypothesised ancestral group likely split into two branches again around . One branch would evolve into cetaceans, possibly beginning about , with the protowhale Pakicetus and other early whale ancestors collectively known as Archaeoceti. This group eventually underwent aquatic adaptation into the completely aquatic cetaceans. The other branch became the anthracotheres, a large family of four-legged beasts, the earliest of which in the late Eocene would have resembled skinny hippos with comparatively smaller, narrower heads. All branches of the anthracotheres, except that which evolved into Hippopotamidae, became extinct during the Pliocene, leaving no descendants. A rough evolutionary lineage of the hippo can thus be traced from Eocene and Oligocene species: from Anthracotherium and Elomeryx to the Miocene species Merycopotamus and Libycosaurus and finally the very latest anthracotheres in the Pliocene. These groups lived across Eurasia and Africa. The discovery of Epirigenys in East Africa, which was likely a descent of Asian anthracotheres and a sister taxon to Hippopotamidae, suggests that hippo ancestors entered Africa from Asia around . An early hippopotamid is the genus Kenyapotamus, which lived in Africa from 15 to . Hippopotamid species would spread across Africa and Eurasia, including the modern pygmy hippo. From 7.5 to , a possible ancestor to the modern hippo, Archaeopotamus, lived in Africa and the Middle East. The oldest records of the genus Hippopotamus date to the Pliocene (5.3–2.6 million years ago). The oldest unambiguous records of the modern H. amphibius date to the Middle Pleistocene, though there are possible Early Pleistocene records. Extinct species Three species of Malagasy hippopotamus became extinct during the Holocene on Madagascar, the last of them within the past 1,000 years. The Malagasy hippos were smaller than the modern hippo, a likely result of the process of insular dwarfism. Fossil evidence indicates many Malagasy hippos were hunted by humans, a factor in their eventual extinction. Isolated individual Malagasy hippos may have survived in remote pockets; in 1976, villagers described a living animal called the kilopilopitsofy, which may have been a Malagasy hippo. Hippopotamus gorgops from the Early Pleistocene to the early Middle Pleistocene of Africa and West Asia grew considerably larger than the living hippopotamus, with an estimated body mass of over . Hippopotamus antiquus ranged throughout Europe, extending as far north as Britain during the Early and Middle Pleistocene epochs, before being replaced by the modern H. amphibius in Europe during the latter part of the Middle Pleistocene. The Pleistocene also saw a number of dwarf species evolve on several Mediterranean islands, including Crete (Hippopotamus creutzburgi), Cyprus (the Cyprus dwarf hippopotamus, Hippopotamus minor), Malta (Hippopotamus melitensis), and Sicily (Hippopotamus pentlandi). Of these, the Cyprus dwarf hippo survived until the end of the Pleistocene or early Holocene. Evidence from the archaeological site Aetokremnos continues to cause debate on whether or not the species was driven to extinction or even encountered by humans. Characteristics The hippopotamus is a megaherbivore and is exceeded in size among land animals only by elephants and some rhinoceros species. The mean adult weight is around for bulls and for cows. Exceptionally large males have been recorded reaching . Male hippos appear to continue growing throughout their lives, while females reach maximum weight at around age 25. It is long, including a tail of about in length and tall at the shoulder, with males and females ranging and tall at the shoulder respectively. The species has a typical head–body length of and an average standing height of at the shoulder. Hippos have barrel-shaped bodies with short tails and legs, and an hourglass-shaped skull with a long snout. Their skeletal structures are graviportal, adapted to carrying their enormous weight, and their dense bones and low centre of gravity allows them to sink and move along the bottom of the water. Hippopotamuses have small legs (relative to other megafauna) because the water in which they live reduces the weight burden. The toes are webbed and the pelvis rests at an angle of 45 degrees. Though chubby-looking, hippos have little fat. The eyes, ears, and nostrils of hippos are placed high on the roof of their skulls. This allows these organs to remain above the surface while the rest of the body is submerged. The nostrils and ears can close when underwater while nictitating membranes cover the eyes. The vocal folds of the hippo are more horizontally positioned, much like baleen whales. Underneath are throat tissues, where vibrations are transmitted to produce underwater calls. The hippo's jaw is powered by huge masseter and digastric muscles which give them large, droopy cheeks. The jaw hinge allows the animal to open its mouth at almost 180°. A folded orbicularis oris muscle allows the hippo to attain an extreme gape without tearing any tissue. On the lower jaw, the incisors and canines grow continuously, the former reaching , while the latter can grow to up to . The lower canines are sharpened through contact with the smaller upper canines. The canines and incisors are used mainly for combat instead of feeding. Hippos rely on their flattened, horny lips to grasp and pull grasses which are then ground by the molars. The hippo is considered to be a pseudoruminant; it has a complex three-chambered stomach, but does not "chew cud". Hippo skin is thick across much of its body with little hair. The animal is mostly purplish-grey or blue-black, but brownish-pink on the underside and around the eyes and ears. Their skin secretes a natural, red-coloured sunscreen substance that is sometimes referred to as "blood sweat" but is neither blood nor sweat. This secretion is initially colourless and turns red-orange within minutes, eventually becoming brown. Two highly acidic pigments have been identified in the secretions; one red hipposudoric acid and one orange norhipposudoric acid, which inhibit the growth of disease-causing bacteria and their light-absorption profile peaks in the ultraviolet range, creating a sunscreen effect. Regardless of diet, all hippos secrete these pigments so food does not appear to be their source; rather, they may be synthesised from precursors such as the amino acid tyrosine. This natural sunscreen cannot prevent the animal's skin from cracking if it stays out of water too long. The testes of the males do not fully descend and a scrotum is not present. In addition, the penis retracts into the body when not erect. The genitals of the female hippos are unusual in that the vagina is ridged and the vulval vestibule has two large, protruding diverticula. Both of these have an unknown function. A hippo's lifespan is typically 40 to 50 years. Donna the Hippo was one of the oldest living hippos in captivity. She lived at the Mesker Park Zoo in Evansville, Indiana, in the US until her death in 2012 at the age of 61. The oldest hippo ever recorded was called Bertha; she had lived in the Manila Zoo in the Philippines since it first opened in 1959. When she died in 2017, her age was estimated to be 65. The oldest living hippopotamus in captivity is Lu the Hippo, from the Ellie Schiller Homosassa Springs Wildlife State Park. As of 2024, he is 64 years old. Distribution and status Hippopotamus amphibius arrived in Europe around 560–460,000 years ago, during the Middle Pleistocene. The distribution of Hippopotamus amphibius in Europe during the Pleistocene was largely confined to Southern Europe, including the Iberian Peninsula, Italy, and Greece, but extended into northwestern Europe, including Great Britain (as far north as North Yorkshire), the Netherlands, and western Germany during interglacial periods, such as the Last Interglacial (130–115,000 years ago). The youngest records of the species in Europe are from the Late Pleistocene of Greece, dating to around 40–30,000 years ago. Archaeological evidence exists of its presence in the Levant, dating to less than 3,000 years ago. The species was common in Egypt's Nile region during antiquity, but it has since been driven out. According to Pliny the Elder, in his time, the best location in Egypt for capturing this animal was in the Saite nome; the animal could still be found along the Damietta branch of the Nile after the Arab Conquest in 639. Reports of the slaughter of the last hippo in Natal Province were made at the end of the 19th century. Hippos are still found in the rivers and lakes of the northern Democratic Republic of the Congo, Uganda, Tanzania, and Kenya, north through to Ethiopia, Somalia, and Sudan, west to The Gambia, and south to South Africa. Genetic evidence suggests common hippos in Africa experienced a marked population expansion during or after the Pleistocene, attributed to an increase in water bodies at the end of the era. These findings have important conservation implications, as hippo populations across the continent are currently threatened by loss of access to fresh water. Hippos are also subject to unregulated hunting and poaching. The species is included in Appendix II of the Convention on International Trade in Endangered Species (CITES) meaning international export/import (including in parts and derivatives) requires CITES documentation to be obtained and presented to border authorities. As of 2017, the IUCN Red List drawn up by the International Union for Conservation of Nature (IUCN) lists the species as vulnerable, with a stable population estimated between 115,000 and 130,000 animals. The hippo population has declined most dramatically in the Democratic Republic of the Congo. By 2005, the population in Virunga National Park had dropped to 800 or 900 from around 29,000 in the mid-1970s. This decline is attributed to the disruptions caused by the Second Congo War. The poachers are believed to be Mai-Mai rebels, underpaid Congolese soldiers, and local militia groups. Reasons for poaching include the belief hippos are harmful to society, as well as financial gain. As of 2016, the Virunga hippo population appears to have increased again, possibly due to better protection from park rangers, who have worked with local fishermen. The sale of hippo meat is illegal, but black-market sales are difficult for Virunga National Park officers to track. Hippo meat is highly valued in some areas of central Africa and the teeth may be used as a replacement for elephant ivory. A population of hippos exists in Colombia, descended from captive individuals that escaped from Pablo Escobar's estate after his death in 1993. Their numbers grew to 100 by the 2020s and ecologists believe the population should be eradicated, as they are breeding rapidly and are an increasing menace to humans and the environment. Attempts to control them include sterilisation and culling. Behaviour and ecology Hippos are semiaquatic and require enough water to immerse in, while being close to grass. They mostly live in freshwater habitat, but can be found in estuaries. They prefer relatively still waters with gently sloping shores, though male hippos may also be found in very small numbers in more rapid waters with rocky slopes. Like most herbivores, hippos will consume a variety of plants if presented with them in captivity, but their diet in nature consists almost entirely of grass, with only minimal consumption of aquatic plants. Hippos spend most of the day in water to stay cool and hydrated. Just before night begins, they leave the water to forage on land. A hippo will travel per night, eating around of grass. By dawn, they are back in the water. Despite being semiaquatic, an adult hippo is not a particularly good swimmer, nor can it float. It rarely enters deep water; when it does, the animal moves by bouncing off the bottom. An adult hippo surfaces every four to six minutes, while young need to breathe every two to three minutes. Hippos move on land by trotting, and limb movements do not change between speeds. They can reach an airborne stage (a stage when all limb are off the ground) when they move fast enough. Hippos are reported to reach but this has not been confirmed. They are incapable of jumping but can walk up steep banks. The hippopotamus sleeps with both hemispheres of the brain resting, as in all land mammals, and usually sleeps on land or in water with the nostrils exposed. Despite this, it may be capable of sleeping while submerged, intermittently surfacing to breathe without waking. They appear to transition between different phases of sleep more quickly than other mammals. Because of their size and their habit of taking the same paths to feed, hippos can have a significant impact on the land across which they walk, keeping the land clear of vegetation and depressing the ground. Over prolonged periods, hippos can divert the paths of swamps and channels. By defecating in the water, the animals also appear to pass on microbes from their gut, affecting the biogeochemical cycle. On occasion, hippos have been filmed eating carrion, usually near the water. There are other reports of meat-eating and even cannibalism and predation. Hippos' stomach anatomy lacks adaptions to carnivory and meat-eating is likely caused by lack of nutrients or just an abnormal behaviour. Social life It is challenging to study the interaction of bulls and cows because hippos are not sexually dimorphic, so cows and young bulls are almost indistinguishable in the field. Hippo pods fluctuate but can contain over 100 hippos. Although they lie close together, adults develop almost no social bonds. Males establish territories in water but not land, and these may range in lakes and in rivers. Territories are abandoned when the water dries up. The bull has breeding access to all the cows in his territory. Younger bachelors are allowed to stay as long as they defer to him. A younger male may challenge the old bull for control of the territory. Within the pods, the hippos tend to segregate by sex and status. Bachelor males lounge near other bachelors, females with other females, and the territorial male is on his own. When hippos emerge from the water to graze, they do so individually. Hippos engage in "muck-spreading" which involves defecating while spinning their tails to distribute the faeces over a greater area. Muck-spreading occurs both on land and in water and its function is not well understood. It is unlikely to serve a territorial function, as the animals only establish territories in the water. They may be used as trails between the water and grazing areas. "Yawning" serves as a threat display. When fighting, bulls use their incisors to block each other's attacks and their large canines as offensive weapons. When hippos become over-populated or a habitat shrinks, bulls sometimes attempt infanticide, but this behaviour is not common under normal conditions. The most common hippo vocalisation is the "wheeze honk", which can travel over long distances in air. This call starts as a high-pitched squeal followed by a deeper, resonant call. The animals can recognise the calls of other individuals. Hippos are more likely to react to the wheeze honks of strangers than to those they are more familiar with. When threatened or alarmed, they produce exhalations, and fighting bulls will bellow loudly. Hippos are recorded to produce clicks underwater which may have echolocative properties. They have the unique ability to hold their heads partially above the water and send out a cry that travels through both water and air; individuals respond both above and below water. Reproduction Cows reach sexual maturity at five to six years of age and have a gestation period of eight months. A study of endocrine systems revealed cows may begin puberty at as early as three or four years. Bulls reach maturity at around 7.5 years. Both conceptions and births are highest during the wet season. Male hippos always have mobile spermatozoa and can breed year-round. After becoming pregnant, a female hippo will typically not begin ovulation again for 17 months. Hippos mate in the water, with the cow remaining under the surface, her head emerging periodically to draw breath. Cows give birth in seclusion and return within 10 to 14 days. Calves are born on land or shallow water weighing on average and at an average length of around . The female lies on her side when nursing, which can occur underwater or on land. The young are carried on their mothers' backs in deep water. Mother hippos are very protective of their young, not allowing others to get too close. One cow was recorded protecting a calf's carcass after it had died. Calves may be temporarily kept in nurseries, guarded by one or more adults, and will play amongst themselves. Like many other large mammals, hippos are described as K-strategists, in this case typically producing just one large, well-developed infant every couple of years (rather than many small, poorly developed young several times per year, as is common among small mammals such as rodents). Calves no longer need to suckle when they are a year old. Interspecies interactions Hippos coexist alongside a variety of large predators in their habitats. Nile crocodiles, lions, and spotted hyenas are known to prey on young hippos. Beyond these, adult hippos are not usually preyed upon by other animals due to their aggression and size. Cases where large lion prides have successfully preyed on adult hippos have been reported, but it is generally rare. Lions occasionally prey on adults at Gorongosa National Park and calves are sometimes taken at Virunga. Crocodiles are frequent targets of hippo aggression, probably because they often inhabit the same riparian habitats; crocodiles may be either aggressively displaced or killed by hippos. In turn, very large Nile crocodiles have been observed preying occasionally on calves, "half-grown" hippos, and possibly also adult female hippos. Groups of crocodiles have also been observed finishing off still-living male hippos that were previously injured in mating battles with other males. Hippos occasionally visit cleaning stations in order to be cleaned of parasites by certain species of fishes. They signal their readiness for this service by opening their mouths wide. This is an example of mutualism, in which the hippo benefits from the cleaning while the fish receive food. Hippo defecation creates allochthonous deposits of organic matter along the river beds. These deposits have an unclear ecological function. A 2015 study concluded hippo dung provides nutrients from terrestrial material for fish and aquatic invertebrates, while a 2018 study found that their dung can be toxic to aquatic life in large quantities, due to absorption of dissolved oxygen in water bodies. The parasitic monogenean flatworm Oculotrema hippopotami infests hippopotamus eyes, mainly the nictitating membrane. It is the only monogenean species (which normally live on fish) documented to live on a mammal. Hippos and humans Cut marks on bones of H. amphibius found at Bolomor Cave, a site in Spain preserving fossils dating from 230,000 to 120,000 years ago, provides evidence for Neanderthal butchery of hippopotamuses. The earliest evidence of modern human interaction with hippos comes from butchery cut marks on hippo bones found at the Bouri Formation and dated to around 160,000 years ago. 4,000–5,000 year art showing hippos being hunted have been found in the Tassili n'Ajjer Mountains of the central Sahara near Djanet. The ancient Egyptians recognised the hippo as ferocious, and representations on the tombs of nobles show humans hunting them. The hippo was also known to the Greeks and Romans. The Greek historian Herodotus described the hippo in The Histories (written circa 440 BC) and the Roman naturalist Pliny the Elder wrote about the hippo in his encyclopedia Naturalis Historia (written circa 77 AD). The Yoruba people called the hippo erinmi, which means "elephant of the water". Some individual hippos have achieved international fame. Huberta became a celebrity during the Great Depression for trekking a great distance across South Africa. Attacks on humans The hippo is considered to be extremely aggressive and has frequently been reported charging and attacking boats. Small boats can easily be capsized by hippos and passengers can be injured or killed by the animals, or drown in the water. In one 2014 case in Niger, a boat was capsized by a hippo and 13 people were killed. Hippos will often raid farm crops if the opportunity arises, and humans may come into conflict with them on these occasions. These encounters can be fatal to either humans or hippos. According to the Ptolemaic historian Manetho, the pharaoh Menes was carried off and then killed by a hippopotamus. In zoos Hippos have long been popular zoo animals. The first record of hippos taken into captivity for display is dated to 3500 BC in Hierakonpolis, Egypt. The first zoo hippo in modern history was Obaysch, who arrived at the London Zoo on 25 May 1850, where he attracted up to 10,000 visitors a day and inspired a popular song, the "Hippopotamus Polka". Hippos generally breed well in captivity; birth rates are lower than in the wild, but this can be attributed to zoos' desire to limit births, since hippos are relatively expensive to maintain. Starting in 2015, the Cincinnati Zoo built a US$73 million exhibit to house three adult hippos, featuring a tank. Modern hippo enclosures also have a complex filtration system for the animal's waste, an underwater viewing area for the visitors, and glass that may be up to thick and capable of holding water under pressures of . Cultural significance In Egyptian mythology, the god Set takes the form of a red hippopotamus and fights Horus for control of the land, but is defeated. The goddess Tawaret is depicted as a pregnant woman with a hippo head, representing fierce maternal love. The Ijaw people of the Niger Delta wore masks of aquatic animals like the hippo when practising their water spirit cults, and hippo ivory was used in the divination rituals of the Yoruba. Hippo masks were also used in Nyau funerary rituals of the Chewa of Southern Africa. According to Robert Baden-Powell, Zulu warriors referred to hippos in war chants. The Behemoth from the Book of Job, 40:15–24 is thought to be based on the hippo. Hippos have been the subjects of various African folktales. According to a San story, when the Creator assigned each animal its place in nature, the hippos wanted to live in the water, but were refused out of fear they might eat all the fish. After begging and pleading, the hippos were finally allowed to live in the water on the condition they would eat grass instead of fish, and fling their dung so it can be inspected for fish bones. In a Ndebele tale, the hippo originally had long, beautiful hair, but it was set on fire by a jealous hare and the hippo had to jump into a nearby pool. The hippo lost most of his hair and was too embarrassed to leave the water. Hippopotamuses were rarely depicted in European art during the Renaissance and Baroque periods, due to less access to specimens by Europeans. One notable exception is Peter Paul Rubens' The Hippopotamus and Crocodile Hunt (1615–1616). Ever since Obaysch inspired the "Hippopotamus Polka", hippos have been popular animals in Western culture for their rotund appearance, which many consider comical. The Disney film Fantasia featured a ballerina hippo dancing to the opera La Gioconda. The film Hugo the Hippo is set in Tanzania and involves the title character trying to escape being slaughtered with the help of local children. The Madagascar films feature a hippo named Gloria. Hippos even inspired a popular board game, Hungry Hungry Hippos. Among the most famous poems about the hippo is "The Hippopotamus" by T. S. Eliot, where he uses the animal to represent the Catholic Church. Hippos are mentioned in the novelty Christmas song "I Want a Hippopotamus for Christmas" that became a hit for child star Gayla Peevey in 1953. They also featured in the popular "The Hippopotamous Song" by Flanders and Swann.	Biology and health sciences	Artiodactyla	null
37789	https://en.wikipedia.org/wiki/Mosquito	Mosquito	Mosquitoes, the Culicidae, are a family of small flies consisting of 3,600 species. The word mosquito (formed by mosca and diminutive -ito) is Spanish and Portuguese for little fly. Mosquitoes have a slender segmented body, one pair of wings, three pairs of long hair-like legs, and specialized, highly elongated, piercing-sucking mouthparts. All mosquitoes drink nectar from flowers; females of some species have in addition adapted to drink blood. The group diversified during the Cretaceous period. Evolutionary biologists view mosquitoes as micropredators, small animals that parasitise larger ones by drinking their blood without immediately killing them. Medical parasitologists view mosquitoes instead as vectors of disease, carrying protozoan parasites or bacterial or viral pathogens from one host to another. The mosquito life cycle consists of four stages: egg, larva, pupa, and adult. Eggs are laid on the water surface; they hatch into motile larvae that feed on aquatic algae and organic material. These larvae are important food sources for many freshwater animals, such as dragonfly nymphs, many fish, and some birds. Adult females of many species have mouthparts adapted to pierce the skin of a host and feed on blood of a wide range of vertebrate hosts, and some invertebrates, primarily other arthropods. Some species only produce eggs after a blood meal. The mosquito's saliva is transferred to the host during the bite, and can cause an itchy rash. In addition, blood-feeding species can ingest pathogens while biting, and transmit them to other hosts. Those species include vectors of parasitic diseases such as malaria and filariasis, and arboviral diseases such as yellow fever and dengue fever. By transmitting diseases, mosquitoes cause the deaths of over 725,000 people each year. Description and life cycle Like all flies, mosquitoes go through four stages in their life cycles: egg, larva, pupa, and adult. The first three stages—egg, larva, and pupa—are largely aquatic, the eggs usually being laid in stagnant water. They hatch to become larvae, which feed, grow, and molt until they change into pupae. The adult mosquito emerges from the mature pupa as it floats at the water surface. Mosquitoes have adult lifespans ranging from as short as a week to around a month. Some species overwinter as adults in diapause. Adult Mosquitoes have one pair of wings, with distinct scales on the surface. Their wings are long and narrow, while the legs are long and thin. The body, usually grey or black, is slender, and typically 3–6 mm long. When at rest, mosquitoes hold their first pair of legs outwards, whereas the somewhat similar Chironomid midges hold these legs forwards. Anopheles mosquitoes can fly for up to four hours continuously at , traveling up to in a night. Males beat their wings between 450 and 600 times per second, driven indirectly by muscles which vibrate the thorax. Mosquitoes are mainly small flies; the largest are in the genus Toxorhynchites, at up to in length and in wingspan. Those in the genus Aedes are much smaller, with a wingspan of . Mosquitoes can develop from egg to adult in hot weather in as few as five days, but it may take up to a month. At dawn or dusk, within days of pupating, males assemble in swarms, mating when females fly in. The female mates only once in her lifetime, attracted by the pheromones emitted by the male. As a species that need blood for the eggs to develop, the female finds a host and drinks a full meal of blood. She then rests for two or three days to digest the meal and allow her eggs to develop. She is then ready to lay the eggs and repeat the cycle of feeding and laying. Females can live for up to three weeks in the wild, depending on temperature, humidity, their ability to obtain a blood meal, and avoiding being killed by their vertebrate hosts. Eggs The eggs of most mosquitoes are laid in stagnant water, which may be a pond, a marsh, a temporary puddle, a water-filled hole in a tree, or the water-trapping leaf axils of a bromeliad. Some lay near the water's edge while others attach their eggs to aquatic plants. A few, like Opifex fuscus, can breed in salt-marshes. Wyeomyia smithii breeds in the pitchers of pitcher plants, its larvae feeding on decaying insects that have drowned there. Oviposition, egg-laying, varies between species. Anopheles females fly over the water, touching down or dapping to place eggs on the surface one at a time; their eggs are roughly cigar-shaped and have floats down their sides. A female can lay 100–200 eggs in her lifetime. Aedes females drop their eggs singly, on damp mud or other surfaces near water; their eggs hatch only when they are flooded. Females in genera such as Culex, Culiseta, and Uranotaenia lay their eggs in floating rafts. Mansonia females in contrast lay their eggs in arrays, attached usually to the under-surfaces of waterlily pads. Clutches of eggs of most mosquito species hatch simultaneously, but Aedes eggs in diapause hatch irregularly over an extended period. Larva The mosquito larva's head has prominent mouth brushes used for feeding, a large thorax with no legs, and a segmented abdomen. It breathes air through a siphon on its abdomen, so must come to the surface frequently. It spends most of its time feeding on algae, bacteria, and other microbes in the water's surface layer. It dives below the surface when disturbed. It swims either by propelling itself with its mouth brushes, or by jerkily wriggling its body. It develops through several stages, or instars, molting each time, after which it metamorphoses into a pupa. Aedes larvae, except when very young, can withstand drying; they go into diapause for several months if their pond dries out. Pupa The head and thorax of the pupa are merged into a cephalothorax, with the abdomen curving around beneath it. The pupa or "tumbler" can swim actively by flipping its abdomen. Like the larva, the pupa of most species must come to the surface frequently to breathe, which they do through a pair of respiratory trumpets on their cephalothoraxes. They do not feed; they pass much of their time hanging from the surface of the water by their respiratory trumpets. If alarmed, they swim downwards by flipping their abdomens in much the same way as the larvae. If undisturbed, they soon float up again. The adult emerges from the pupa at the surface of the water and flies off. Feeding by adults Diet Both male and female mosquitoes feed on nectar, aphid honeydew, and plant juices, but in many species the females are also blood-sucking ectoparasites. In some of those species, a blood meal is essential for egg production; in others, it just enables the female to lay more eggs. Both plant materials and blood are useful sources of energy in the form of sugars. Blood supplies more concentrated nutrients, such as lipids, but the main function of blood meals is to obtain proteins for egg production. Mosquitoes like Toxorhynchites reproduce autogenously, not needing blood meals. Disease vector mosquitoes like Anopheles and Aedes are anautogenous, requiring blood to lay eggs. Many Culex species are partially anautogenous, needing blood only for their second and subsequent clutches of eggs. Host animals Blood-sucking mosquitoes favour particular host species, though they are less selective when food is short. Different mosquito species favor amphibians, reptiles including snakes, birds, and mammals. For example, Culiseta melanura sucks the blood of passerine birds, but as mosquito numbers rise they attack mammals including horses and humans, causing epidemics of Eastern equine encephalitis virus in North America. Loss of blood from many bites can add up to a large volume, occasionally causing the death of livestock as large as cattle and horses. Malaria-transmitting mosquitoes seek out caterpillars and feed on their haemolymph, impeding their development. Finding hosts Most mosquito species are crepuscular, feeding at dawn or dusk, and resting in a cool place through the heat of the day. Some species, such as the Asian tiger mosquito, are known to fly and feed during daytime. Female mosquitoes hunt for hosts by smelling substances such as carbon dioxide (CO2) and 1-octen-3-ol (mushroom alcohol, found in exhaled breath) produced from the host, and through visual recognition. The semiochemical that most strongly attracts Culex quinquefasciatus is nonanal. Another attractant is sulcatone. A large part of the mosquito's sense of smell, or olfactory system, is devoted to sniffing out blood sources. Of 72 types of odor receptors on its antennae, at least 27 are tuned to detect chemicals found in perspiration. In Aedes, the search for a host takes place in two phases. First, the mosquito flies about until it detects a host's odorants; then it flies towards them, using the concentration of odorants as its guide. Mosquitoes prefer to feed on people with type O blood, an abundance of skin bacteria, high body heat, and pregnant women. Individuals' attractiveness to mosquitoes has a heritable, genetically controlled component. The multitude of characteristics in a host observed by the mosquito allows it to select a host to feed on. This occurs when a mosquito notes the presence of CO2, as it then activates odour and visual search behaviours that it otherwise would not use. In terms of a mosquito’s olfactory system, chemical analysis has revealed that people who are highly attractive to mosquitoes produce significantly more carboxylic acids. A human's unique body odour indicates that the target is actually a human host rather than some other living warm-blooded animal (as the presence of CO2 shows). Body odour, composed of volatile organic compounds emitted from the skin of humans, is the most important cue used by mosquitoes. Variation in skin odour is caused by body weight, hormones, genetic factors, and metabolic or genetic disorders. Infections such as malaria can influence an individual’s body odour. People infected by malaria produce relatively large amounts of Plasmodium-induced aldehydes in the skin, creating large cues for mosquitoes as it increases the attractiveness of an odour blend, imitating a "healthy" human odour. Infected individuals produce larger amounts of aldehydes heptanal, octanal, and nonanal. These compounds are detected by mosquito antennae. Thus, people infected with malaria are more prone to mosquito biting. Contributing to a mosquito's ability to activate search behaviours, a mosquito's visual search system includes sensitivity to wavelengths from different colours. Mosquitoes are attracted to longer wavelengths, correlated to the colours of red and orange as seen by humans, and range through the spectrum of human skin tones. In addition, they have a strong attraction to dark, high-contrast objects, because of how longer wavelengths are perceived against a lighter-coloured background. Different species of mosquitoes have evolved different methods of identifying target hosts. Study of a domestic form and an animal-biting form of the mosquito Aedes aegypti showed that the evolution of preference for human odour is linked to increases in the expression of the olfactory receptor AaegOr4. This recognises a compound present at high levels in human odour called sulcatone. However, the malaria mosquito Anopheles gambiae also has OR4 genes strongly activated by sulcatone, yet none of them are closely related to AaegOr4, suggesting that the two species have evolved to specialise in biting humans independently. Mouthparts Female mosquito mouthparts are highly adapted to piercing skin and sucking blood. Males only drink sugary fluids, and have less specialized mouthparts. Externally, the most obvious feeding structure of the mosquito is the proboscis, composed of the labium, U-shaped in section like a rain gutter, which sheaths a bundle (fascicle) of six piercing mouthparts or stylets. These are two mandibles, two maxillae, the hypopharynx, and the labrum. The labium bends back into a bow when the mosquito begins to bite, staying in contact with the skin and guiding the stylets downwards. The extremely sharp tips of the labrum and maxillae are moved backwards and forwards to saw their way into the skin, with just one thousandth of the force that would be needed to penetrate the skin with a needle, resulting in a painless insertion. Saliva Mosquito saliva contains enzymes that aid in sugar feeding, and antimicrobial agents that control bacterial growth in the sugar meal. For a mosquito to obtain a blood meal, it must circumvent its vertebrate host's physiological responses. Mosquito saliva blocks the host's hemostasis system, with proteins that reduce vascular constriction, blood clotting, and platelet aggregation, to ensure the blood keeps flowing. It modulates the host's immune response via a mixture of proteins which lower angiogenesis and immunity; create inflammation; suppress tumor necrosis factor release from activated mast cells; suppress interleukin (IL)-2 and IFN-γ production; suppress T cell populations; decrease expression of interferon−α/β, making virus infections more severe; increase natural killer T cells in the blood; and decrease cytokine production. Egg development and blood digestion Females of many blood-feeding species need a blood meal to begin the process of egg development. A sufficiently large blood meal triggers a hormonal cascade that leads to egg development. Upon completion of feeding, the mosquito withdraws her proboscis, and as the gut fills up, the stomach lining secretes a peritrophic membrane that surrounds the blood. This keeps the blood separate from anything else in the stomach. Like many Hemiptera that survive on dilute liquid diets, many adult mosquitoes excrete surplus liquid even when feeding. This permits females to accumulate a full meal of nutrient solids. The blood meal is digested over a period of several days. Once blood is in the stomach, the midgut synthesizes protease enzymes, primarily trypsin assisted by aminopeptidase, that hydrolyze the blood proteins into free amino acids. These are used in the synthesis of vitellogenin, which in turn is made into egg yolk protein. Distribution Cosmopolitan Mosquitoes have a cosmopolitan distribution, occurring in every land region except Antarctica and a few islands with polar or subpolar climates, such as Iceland, which is essentially free of mosquitoes. This absence is probably caused by Iceland's climate. Its weather is unpredictable, freezing but often warming suddenly in mid-winter, making mosquitoes emerge from pupae in diapause, and then freezing again before they can complete their life cycle. Eggs of temperate zone mosquitoes are more tolerant of cold than the eggs of species indigenous to warmer regions. Many can tolerate subzero temperatures, while adults of some species can survive winter by sheltering in microhabitats such as buildings or hollow trees. In warm and humid tropical regions, some mosquito species are active for the entire year, but in temperate and cold regions they hibernate or enter diapause. Arctic or subarctic mosquitoes, like some other arctic midges in families such as Simuliidae and Ceratopogonidae may be active for only a few weeks annually as melt-water pools form on the permafrost. During that time, though, they emerge in huge numbers in some regions; a swarm may take up to 300 ml of blood per day from each animal in a caribou herd. Effect of climate change For a mosquito to transmit disease, there must be favorable seasonal conditions, primarily humidity, temperature, and precipitation. El Niño affects the location and number of outbreaks in East Africa, Latin America, Southeast Asia and India. Climate change impacts the seasonal factors and in turn the dispersal of mosquitoes. Climate models can use historic data to recreate past outbreaks and to predict the risk of vector-borne disease, based on an area's forecasted climate. Mosquito-borne diseases have long been most prevalent in East Africa, Latin America, Southeast Asia, and India. An emergence in Europe was observed early in the 21st century. It is predicted that by 2030, the climate of southern Great Britain will be suitable for transmission of Plasmodium vivax malaria by Anopheles mosquitoes for two months of the year, and that by 2080, the same will be true for southern Scotland. Dengue fever, too, is spreading northwards with climate change. The vector, the Asian tiger mosquito Aedes albopictus, has by 2023 established across southern Europe and as far north as much of northern France, Belgium, Holland, and both Kent and West London in England. Ecology Predators and parasites Mosquito larvae are among the commonest animals in ponds, and they form an important food source for freshwater predators. Among the many aquatic insects that catch mosquito larvae are dragonfly and damselfly nymphs, whirligig beetles, and water striders. Vertebrate predators include fish such as catfish and the mosquitofish, amphibians including the spadefoot toad and the giant tree frog, freshwater turtles such as the red-eared slider, and birds such as ducks. Emerging adults are consumed at the pond surface by predatory flies including Empididae and Dolichopodidae, and by spiders. Flying adults are captured by dragonflies and damselflies, by birds such as swifts and swallows, and by vertebrates including bats. Mosquitoes are parasitised by hydrachnid mites, ciliates such as Glaucoma, microsporidians such as Thelania, and fungi including species of Saprolegniaceae and Entomophthoraceae. Pollination Several flowers including members of the Asteraceae, Rosaceae and Orchidaceae are pollinated by mosquitoes, which visit to obtain sugar-rich nectar. They are attracted to flowers by a range of semiochemicals such as alcohols, aldehydes, ketones, and terpenes. Mosquitoes have visited and pollinated flowers since the Cretaceous period. It is possible that plant-sucking exapted mosquitoes to blood-sucking. Parasitism Ecologically, blood-feeding mosquitoes are micropredators, small animals that feed on larger animals without immediately killing them. Evolutionary biologists see this as a form of parasitism; in Edward O. Wilson's phrase "Parasites ... are predators that eat prey in units of less than one." Micropredation is one of six major evolutionarily stable strategies within parasitism. It is distinguished by leaving the host still able to reproduce, unlike the activity of parasitic castrators or parasitoids; and having multiple hosts, unlike conventional parasites. From this perspective, mosquitoes are ectoparasites, feeding on blood from the outside of their hosts, using their piercing mouthparts, rather than entering their bodies. Unlike some other ectoparasites such as fleas and lice, mosquitoes do not remain constantly on the body of the host, but visit only to feed. Evolution Fossil record A 2023 study suggested that Libanoculex intermedius found in Lebanese amber, dating to the Barremian age of the Early Cretaceous, around 125 million years ago was the oldest known mosquito. However its identification as a mosquito is disputed, with other authors considering it to be a chaoborid fly instead. Three other unambiguous species of Cretaceous mosquito are known. Burmaculex antiquus and Priscoculex burmanicus are known from Burmese amber from Myanmar, which dates to the earliest part of the Cenomanian age of the Late Cretaceous, around 99 million years ago. Paleoculicis minutus, is known from Canadian amber from Alberta, Canada, which dates to the Campanian age of the Late Cretaceous, around 79 million years ago. P. burmanicus has been assigned to the Anophelinae, indicating that the split between this subfamily and the Culicinae took place over 99 million years ago. Molecular estimates suggest that this split occurred 197.5 million years ago, during the Early Jurassic, but that major diversification did not take place until the Cretaceous. Taxonomy Over 3,600 species of mosquitoes in 112 genera have been described. They are traditionally divided into two subfamilies, the Anophelinae and the Culicinae, which carry different diseases. Roughly speaking, protozoal diseases like malaria are transmitted by anophelines, while viral diseases such as yellow fever and dengue fever are transmitted by culicines. The name Culicidae was introduced by the German entomologist Johann Wilhelm Meigen in his seven-volume classification published in 1818–1838. Mosquito taxonomy was advanced in 1901 when the English entomologist Frederick Vincent Theobald published his 5-volume monograph on the Culicidae. He had been provided with mosquito specimens sent in to the British Museum (Natural History) from around the world, on the 1898 instruction of the Secretary of State for the Colonies, Joseph Chamberlain, who had written that "in view of the possible connection of Malaria with mosquitoes, it is desirable to obtain exact knowledge of the different species of mosquitoes and allied insects in the various tropical colonies. I will therefore ask you ... to have collections made of the winged insects in the Colony which bite men or animals." Phylogeny External Mosquitoes are members of a family of the true flies (order Diptera): the Culicidae (from the Latin , genitive , meaning "midge" or "gnat"). They are members of the infraorder Culicomorpha and superfamily Culicoidea. The phylogenetic tree is based on the FLYTREE project. Internal The two subfamilies of mosquitoes are Anophelinae, containing three genera and approximately 430 species, and Culicinae, which contains 11 tribes, 108 genera and 3,046 species. Kyanne Reidenbach and colleagues analysed mosquito phylogenetics in 2009, using both nuclear DNA and morphology of 26 species. They note that Anophelinae is confirmed to be rather basal, but that the deeper parts of the tree are not well resolved. Interactions with humans Vectors of disease Mosquitoes are vectors for many disease-causing microorganisms including bacteria, viruses, and protozoan parasites. Nearly 700 million people acquire a mosquito-borne illness each year, resulting in over 725,000 deaths. Common mosquito-borne viral diseases include yellow fever and dengue fever transmitted mostly by Aedes aegypti. Parasitic diseases transmitted by mosquitoes include malaria and lymphatic filariasis. The Plasmodium parasites that cause malaria are carried by female Anopheles mosquitoes. Lymphatic filariasis, the main cause of elephantiasis, is spread by a wide variety of mosquitoes. A bacterial disease spread by Culex and Culiseta mosquitoes is tularemia. Control Many measures have been tried for mosquito control, including the elimination of breeding places, exclusion via window screens and mosquito nets, biological control with parasites such as fungi and nematodes, or predators such as fish, copepods, dragonfly nymphs and adults, and some species of lizard and gecko. Another approach is to introduce large numbers of sterile males. Genetic modification methods including cytoplasmic incompatibility, chromosomal translocations, sex distortion and gene replacement, solutions seen as inexpensive and not subject to vector resistance, have been explored. Control of disease-carrying mosquitoes using gene drives has been proposed. Repellents Insect repellents are applied on skin and give short-term protection against mosquito bites. The chemical DEET repels some mosquitoes and other insects. Some CDC-recommended repellents are picaridin, eucalyptus oil (PMD), and ethyl butylacetylaminopropionate (IR3535). Pyrethrum (from Chrysanthemum species, particularly C. cinerariifolium and C. coccineum) is an effective plant-based repellent. Electronic insect repellent devices that produce ultrasounds intended to keep away insects (and mosquitoes) are marketed. No EPA or university study has shown that these devices prevent humans from being bitten by a mosquito. Bites Mosquito bites lead to a variety of skin reactions and more seriously to mosquito bite allergies. Such hypersensitivity to mosquito bites is an excessive reaction to mosquito saliva proteins. Numerous species of mosquito can trigger such reactions, including Aedes aegypti, A. vexans, A. albopictus, Anopheles sinensis, Culex pipiens, Aedes communis, Anopheles stephensi, C. quinquefasciatus, C. tritaeniorhynchus, and Ochlerotatus triseriatus. Cross-reactivity between salivary proteins of different mosquitoes implies that allergic responses may be caused by virtually any mosquito species. Treatment can be with anti-itch medications, including some taken orally, such as diphenhydramine, or applied to the skin like antihistamines or corticosteroids such as hydrocortisone. Aqueous ammonia (3.6%) also provides relief. Both topical heat and cold may be useful as treatments. In human culture Greek mythology Ancient Greek beast fables including "The Elephant and the Mosquito" and "The Bull and the Mosquito", with the general moral that the large beast does not even notice the small one, derive ultimately from Mesopotamia. Origin myths The peoples of Siberia have origin myths surrounding the mosquito. One Ostiak myth tells of a man-eating giant, Punegusse, who is killed by a hero but will not stay dead. The hero eventually burns the giant, but the ashes of the fire become mosquitoes that continue to plague mankind. Other myths from the Yakuts, Goldes (Nanai people), and Samoyed have the insect arising from the ashes or fragments of some giant creature or demon. Similar tales found in Native North American myth, with the mosquito arising from the ashes of a man-eater, suggest a common origin. The Tatars of the Altai had a variant of the same myth, involving the fragments of the dead giant, Andalma-Muus, becoming mosquitoes and other insects. Lafcadio Hearn tells that in Japan, mosquitoes are seen as reincarnations of the dead, condemned by the errors of their former lives to the condition of Jiki-ketsu-gaki, or "blood-drinking pretas". Modern era Winsor McCay's 1912 film How a Mosquito Operates was one of the earliest works of animation. It has been described as far ahead of its time in technical quality. It depicts a giant mosquito tormenting a sleeping man. Twelve ships of the Royal Navy have borne the name HMS Mosquito or the archaic form of the name, HMS Musquito. The de Havilland Mosquito was a high-speed aircraft manufactured between 1940 and 1950, and used in many roles. The Russian city of Berezniki annually celebrates its mosquitoes from the 17th of July to the 20th in a "most delicious girl" competition. In the competition, the girls stand for 20 minutes in their shorts and vests, and the one who receives the most bites wins.	Biology and health sciences	Flies (Diptera)	null
37796	https://en.wikipedia.org/wiki/West%20Nile%20fever	West Nile fever	West Nile fever is an infection by the West Nile virus, which is typically spread by mosquitoes. In about 80% of infections people have few or no symptoms. About 20% of people develop a fever, headache, vomiting, or a rash. In less than 1% of people, encephalitis or meningitis occurs, with associated neck stiffness, confusion, or seizures. Recovery may take weeks to months. The risk of death among those in whom the nervous system is affected is about 10 percent. West Nile virus (WNV) is usually spread by mosquitoes that become infected when they feed on infected birds, which often carry the disease. Rarely the virus is spread through blood transfusions, organ transplants, or from mother to baby during pregnancy, delivery, or breastfeeding, but it otherwise does not spread directly between people. Risks for severe disease include being over 60 years old and having other health problems. Diagnosis is typically based on symptoms and blood tests. There is no human vaccine. The best way to reduce the risk of infection is to avoid mosquito bites. Mosquito populations may be reduced by eliminating standing pools of water, such as in old tires, buckets, gutters, and swimming pools. When mosquitoes cannot be avoided, mosquito repellent, window screens, and mosquito nets reduce the likelihood of being bitten. There is no specific treatment for the disease; pain medications may reduce symptoms. The virus was discovered in Uganda in 1937, and was first detected in North America in 1999. WNV has occurred in Europe, Africa, Asia, Australia, and North America. In the United States thousands of cases are reported a year, with most occurring in August and September. It can occur in outbreaks of disease. Severe disease may also occur in horses, for which a vaccine is available. A surveillance system in birds is useful for early detection of a potential human outbreak. Signs and symptoms About 80% of those infected with West Nile virus (WNV) show no symptoms and go unreported. About 20% of infected people develop symptoms. These vary in severity, and begin 3 to 14 days after being bitten. Most people with mild symptoms of WNV recover completely, though fatigue and weakness may last for weeks or months. Symptoms may range from mild, such as fever, to severe, such as paralysis and meningitis. A severe infection can last weeks and can, rarely, cause permanent brain damage. Death may ensue if the central nervous system is affected. Medical conditions such as cancer and diabetes, and age over 60 years, increase the risk of developing severe symptoms. Headache can be a prominent symptom of WNV fever, meningitis, encephalitis, meningoencephalitis, and it may or may not be present in poliomyelitis-like syndrome. Thus, headache is not a useful indicator of neuroinvasive disease. West Nile fever (WNF), which occurs in 20 percent of cases, is a febrile syndrome that causes flu-like symptoms. Most characterizations of WNF describe it as a mild, acute syndrome lasting 3 to 6 days after symptom onset. Systematic follow-up studies of patients with WNF have not been done, so this information is largely anecdotal. Possible symptoms include high fever, headache, chills, excessive sweating, weakness, fatigue, swollen lymph nodes, drowsiness, pain in the joints and flu-like symptoms. There may be gastrointestinal symptoms including nausea, vomiting, loss of appetite, and diarrhea. Fewer than one-third of patients develop a rash. West Nile neuroinvasive disease (WNND), which occurs in less than 1 percent of cases, is when the virus infects the central nervous system resulting in meningitis, encephalitis, meningoencephalitis or a poliomyelitis-like syndrome. Many patients with WNND have normal neuroimaging studies, although abnormalities may be present in various cerebral areas including the basal ganglia, thalamus, cerebellum, and brainstem. West Nile virus encephalitis (WNE) is the most common neuroinvasive manifestation of WNND. WNE presents with similar symptoms to other viral encephalitis with fever, headaches, and altered mental status. A prominent finding in WNE is muscular weakness (30 to 50 percent of patients with encephalitis), often with lower motor neuron symptoms, flaccid paralysis, and hyporeflexia with no sensory abnormalities. West Nile meningitis (WNM) usually involves fever, headache, stiff neck and pleocytosis, an increase of white blood cells in cerebrospinal fluid. Changes in consciousness are not usually seen and are mild when present. West Nile meningoencephalitis is inflammation of both the brain (encephalitis) and meninges (meningitis). West Nile poliomyelitis (WNP), an acute flaccid paralysis syndrome associated with WNV infection, is less common than WNM or WNE. This syndrome is generally characterized by the acute onset of asymmetric limb weakness or paralysis in the absence of sensory loss. Pain sometimes precedes the paralysis. The paralysis can occur in the absence of fever, headache, or other common symptoms associated with WNV infection. Involvement of respiratory muscles, leading to acute respiratory failure, sometimes occurs. West-Nile reversible paralysis, Like WNP, the weakness or paralysis is asymmetric. Reported cases have been noted to have an initial preservation of deep tendon reflexes, which is not expected for a pure anterior horn involvement. Disconnect of upper motor neuron influences on the anterior horn cells possibly by myelitis or glutamate excitotoxicity have been suggested as mechanisms. The prognosis for recovery is excellent. Nonneurologic complications of WNV infection that may rarely occur include fulminant hepatitis, pancreatitis, myocarditis, rhabdomyolysis, orchitis, nephritis, optic neuritis and cardiac dysrhythmias and hemorrhagic fever with coagulopathy. Chorioretinitis may also be more common than previously thought. Skin manifestations, specifically rashes, are common; however, there are few detailed descriptions in case reports, and few images are available. Punctate erythematous, macular, and papular eruptions, most pronounced on the extremities have been observed in WNV cases and in some cases histopathologic findings have shown a sparse superficial perivascular lymphocytic infiltrate, a manifestation commonly seen in viral exanthems. A literature review provides support that this punctate rash is a common cutaneous presentation of WNV infection. Cause Virology WNV is one of the Japanese encephalitis antigenic serocomplex of viruses. Image reconstructions and cryoelectron microscopy reveal a 45–50 nm virion covered with a relatively smooth protein surface. This structure is similar to the dengue fever virus; both belong to the genus Flavivirus within the family Flaviviridae. The genetic material of WNV is a positive-sense, single strand of RNA, which is between 11,000 and 12,000 nucleotides long; these genes encode seven nonstructural proteins and three structural proteins. The RNA strand is held within a nucleocapsid formed from 12-kDa protein blocks; the capsid is contained within a host-derived membrane altered by two viral membrane proteins. West Nile virus has been seen to replicate faster and spread more easily to birds at higher temperatures; one of several ways climate change could affect the epidemiology of this disease. Transmission The prime method of spread of the West Nile virus (WNV) is the female mosquito. In Europe, cats were identified as being hosts for West Nile virus. The important mosquito vectors vary according to area; in the United States, Culex pipiens (Eastern United States, and urban and residential areas of the United States north of 36–39°N), Culex tarsalis (Midwest and West), and Culex quinquefasciatus (Southeast) are the main vector species. In Europe, Culex pipiens is the principal vector. The mosquito species that are most frequently infected with WNV feed primarily on birds. Different species of mosquitos take a blood meal from different types of vertebrate hosts, Mosquitoes show further selectivity, exhibiting preference for different species of birds. In the United States, WNV mosquito vectors feed preferentially on members of the Corvidae and thrush family. Among the preferred species within these families are the American crow, a corvid, and the American robin (Turdus migratorius). Some species of birds develop sufficient viral levels (>~104.2 log PFU/ml;) after being infected to transmit the infection to biting mosquitoes that in turn go on to infect other birds. In birds that die from WNV, death usually occurs after 4 to 6 days. In mammals and several species of birds, the virus does not multiply as readily and so does not develop high viremia during infection. Mosquitoes biting such hosts are not believed to ingest sufficient virus to become infected, making them so-called dead-end hosts. As a result of the differential infectiousness of hosts, the feeding patterns of mosquitoes play an important role in WNV transmission, and they are partly genetically controlled, even within a species. Direct human-to-human transmission initially was believed to be caused only by occupational exposure, such as in a laboratory setting, or conjunctival exposure to infected blood. The US outbreak identified additional transmission methods through blood transfusion, organ transplant, intrauterine exposure, and breast feeding. Since 2003, blood banks in the United States routinely screen for the virus among their donors. As a precautionary measure, the UK's National Blood Service initially ran a test for this disease in donors who donate within 28 days of a visit to the United States, Canada, or the northeastern provinces of Italy, and the Scottish National Blood Transfusion Service asks prospective donors to wait 28 days after returning from North America or the northeastern provinces of Italy before donating. There also have been reports of possible transmission of the virus from mother to child during pregnancy or breastfeeding or exposure to the virus in a lab, but these are rare cases and not conclusively confirmed. Recently, the potential for mosquito saliva to affect the course of WNV disease was demonstrated. Mosquitoes inoculate their saliva into the skin while obtaining blood. Mosquito saliva is a pharmacological cocktail of secreted molecules, principally proteins, that can affect vascular constriction, blood coagulation, platelet aggregation, inflammation, and immunity. It clearly alters the immune response in a manner that may be advantageous to a virus. Studies have shown it can specifically modulate the immune response during early virus infection, and mosquito feeding can exacerbate WNV infection, leading to higher viremia and more severe forms of disease. Vertical transmission Vertical transmission, the transmission of a viral or bacterial disease from the female of the species to her offspring, has been observed in various West Nile virus studies, amongst different species of mosquitoes in both the laboratory and in nature. Mosquito progeny infected vertically in autumn may potentially serve as a mechanism for WNV to overwinter and initiate enzootic horizontal transmission the following spring, although it likely plays little role in transmission in the summer and fall. Risk factors Risk factors independently associated with developing a clinical infection with WNV include a suppressed immune system and a patient history of organ transplantation. For neuroinvasive disease the additional risk factors include older age (>50+), male sex, hypertension, and diabetes mellitus. A genetic factor also appears to increase susceptibility to West Nile disease. A mutation of the gene CCR5 gives some protection against HIV but leads to more serious complications of WNV infection. Carriers of two mutated copies of CCR5 made up 4.0 to 4.5% of a sample of people with West Nile disease, while the incidence of the gene in the general population is only 1.0%. The most at risk occupations in the U.S. are outdoor workers, for example farmers, loggers, landscapers/groundskeepers, construction workers, painters, summer camp workers and pavers. Two reports of accidental exposure by laboratory personnel working with infected fluids or tissues have been received. While this appears to be a rare occurrence, it highlights the need for proper handling of infected materials. The World Health Organization states that there are no known cases of health care workers acquiring the virus from infected patients when the appropriate infection control precautions are observed. Diagnosis Preliminary diagnosis is often based on the patient's clinical symptoms, places and dates of travel (if patient is from a nonendemic country or area), activities, and epidemiologic history of the location where infection occurred. A recent history of mosquito bites and an acute febrile illness associated with neurologic signs and symptoms should cause clinical suspicion of WNV. Diagnosis of West Nile virus infections is generally accomplished by serologic testing of blood serum or cerebrospinal fluid (CSF), which is obtained via a lumbar puncture. Initial screening could be done using the ELISA technique detecting immunoglobulins in the sera of the tested individuals. Typical findings of WNV infection include lymphocytic pleocytosis, elevated protein level, reference glucose and lactic acid levels, and no erythrocytes. Definitive diagnosis of WNV is obtained through detection of virus-specific antibody IgM and neutralizing antibodies. Cases of West Nile virus meningitis and encephalitis that have been serologically confirmed produce similar degrees of CSF pleocytosis and are often associated with substantial CSF neutrophilia. Specimens collected within eight days following onset of illness may not test positive for West Nile IgM, and testing should be repeated. A positive test for West Nile IgG in the absence of a positive West Nile IgM is indicative of a previous flavivirus infection and is not by itself evidence of an acute West Nile virus infection. If cases of suspected West Nile virus infection, sera should be collected on both the acute and convalescent phases of the illness. Convalescent specimens should be collected 2–3 weeks after acute specimens. It is common in serologic testing for cross-reactions to occur among flaviviruses such as dengue virus (DENV) and tick-borne encephalitis virus; this necessitates caution when evaluating serologic results of flaviviral infections. Four FDA-cleared WNV IgM ELISA kits are commercially available from different manufacturers in the U.S., each of these kits is indicated for use on serum to aid in the presumptive laboratory diagnosis of WNV infection in patients with clinical symptoms of meningitis or encephalitis. Positive WNV test results obtained via use of these kits should be confirmed by additional testing at a state health department laboratory or CDC. In fatal cases, nucleic acid amplification, histopathology with immunohistochemistry, and virus culture of autopsy tissues can also be useful. Only a few state laboratories or other specialized laboratories, including those at CDC, are capable of doing this specialized testing. Differential diagnosis A number of various diseases may present with symptoms similar to those caused by a clinical West Nile virus infection. Those causing neuroinvasive disease symptoms include the enterovirus infection and bacterial meningitis. Accounting for differential diagnoses is a crucial step in the definitive diagnosis of WNV infection. Consideration of a differential diagnosis is required when a patient presents with unexplained febrile illness, extreme headache, encephalitis or meningitis. Diagnostic and serologic laboratory testing using polymerase chain reaction (PCR) testing and viral culture of CSF to identify the specific pathogen causing the symptoms, is the only currently available means of differentiating between causes of encephalitis and meningitis. Prevention Many of the guidelines for preventing occupational West Nile virus exposure are common to all mosquito-borne diseases. Public health measures include taking steps to reduce mosquito populations. Personal recommendations are to reduce the likelihood of being bitten. General measures to avoid bites include: Using insect repellent on exposed skin to repel mosquitoes. Repellents include products containing DEET and picaridin. DEET concentrations of 30% to 50% are effective for several hours. Picaridin, available at 7% and 15% concentrations, needs more frequent application. DEET formulations as high as 30% are recommended for children over two months of age. The CDC also recommends the use of: IR3535, oil of lemon eucalyptus, para-menthane-diol, or 2-undecanone. Protect infants less than two months of age by using a carrier draped with mosquito netting with an elastic edge for a tight fit. When using sunscreen, apply sunscreen first and then repellent. Repellent should be washed off at the end of the day before going to bed. Wear long-sleeve shirts, which should be tucked in, long trousers, socks, and hats to cover exposed skin (although most fabrics do not totally protect against bites). Insect repellents should be applied over top of protective clothing for greater protection. Do not apply insect repellents underneath clothing. Repellents containing permethrin (e.g., Permanone) or other insect repellents may be applied to clothing, shoes, tents, mosquito nets, and other gear. (Permethrin is not suitable for use directly on skin.) Most repellent is generally removed from clothing and gear by a single washing, but permethrin-treated clothing is effective for up to five washings. Most mosquitoes that transmit disease are most active at dawn and in the evening dusk. A notable exception is the Asian tiger mosquito, which is a daytime feeder and is more apt to be found in, or on the periphery of, shaded areas with heavy vegetation. They are now widespread in the United States, and in Florida they have been found in all 67 counties. In an at-risk area, staying in air-conditioned or well-screened room, or sleeping under an insecticide-treated bed net is recommended. Bed nets should be tucked under mattresses, and can be sprayed with a repellent if not already treated with an insecticide. Monitoring and control West Nile virus can be sampled from the environment by the pooling of trapped mosquitoes via ovitraps, carbon dioxide-baited light traps, and gravid traps, testing blood samples drawn from wild birds, dogs, and sentinel monkeys, and testing brains of dead birds found by various animal control agencies and the public. Testing of the mosquito samples requires the use of reverse-transcriptase PCR (RT-PCR) to directly amplify and show the presence of virus in the submitted samples. When using the blood sera of wild birds and sentinel chickens, samples must be tested for the presence of WNV antibodies by use of immunohistochemistry (IHC) or enzyme-linked immunosorbent assay (ELISA). Dead birds, after necropsy, or their oral swab samples collected on specific RNA-preserving filter paper card, can have their virus presence tested by either RT-PCR or IHC, where virus shows up as brown-stained tissue because of a substrate-enzyme reaction. West Nile control is achieved through mosquito control, by elimination of mosquito breeding sites such as abandoned pools, applying larvacide to active breeding areas, and targeting the adult population via lethal ovitraps and aerial spraying of pesticides. With aerial pesticides, there is a rising need to develop new versions as pesticide resistance among mosquitoes can occur. Environmentalists have condemned attempts to control the transmitting mosquitoes by spraying pesticide, saying the detrimental health effects of spraying outweigh the relatively few lives that may be saved, and more environmentally friendly ways of controlling mosquitoes are available. They also question the effectiveness of insecticide spraying, as they believe mosquitoes that are resting or flying above the level of spraying will not be killed; the most common vector in the northeastern United States, Culex pipiens, is a canopy feeder. Treatment No specific treatment is available for WNV infection. Most people recover without treatment. In mild cases, over-the-counter pain relievers can help ease mild headaches and muscle aches in adults. In severe cases supportive care is provided, often in hospital, with intravenous fluids, pain medication, respiratory support, and prevention of secondary infections. Prognosis While the general prognosis is favorable, current studies indicate that West Nile Fever can often be more severe than previously recognized, with studies of various recent outbreaks indicating that it may take as long as 60 to 90 days to recover. Patients with milder WNF are just as likely as those with more severe manifestations of neuroinvasive disease to experience multiple somatic complaints such as tremor, and dysfunction in motor skills and executive functions for over a year. People with milder symptoms are just as likely as people with more severe symptoms to experience adverse outcomes. Recovery is marked by a long convalescence with fatigue. One study found that neuroinvasive WNV infection was associated with an increased risk for subsequent kidney disease. Epidemiology WNV was first isolated from a feverish 37-year-old woman at Omogo in the West Nile District of Uganda in 1937 during research on yellow fever virus. A series of serosurveys in 1939 in central Africa found anti-WNV positive results ranging from 1.4% (Congo) to 46.4% (White Nile region, Sudan). It was subsequently identified in Egypt (1942) and India (1953), a 1950 serosurvey in Egypt found 90% of those over 40 years in age had WNV antibodies. The ecology was characterized in 1953 with studies in Egypt and Israel. The virus became recognized as a cause of severe human meningoencephalitis in elderly patients during an outbreak in Israel in 1957. The disease was first noted in horses in Egypt and France in the early 1960s and found to be widespread in southern Europe, southwest Asia and Australia. The first appearance of WNV in the Western Hemisphere was in 1999 with encephalitis reported in humans, dogs, cats, and horses, and the subsequent spread in the United States may be an important milestone in the evolving history of this virus. The American outbreak began in College Point, Queens in New York City and was later spread to the neighboring states of New Jersey and Connecticut. The virus is believed to have entered in an infected bird or mosquito, although there is no clear evidence. West Nile virus is now endemic in Africa, Europe, the Middle East, west and central Asia, Oceania (subtype Kunjin), and most recently, North America and is spreading into Central and South America. Outbreaks of West Nile virus encephalitis in humans have occurred in Algeria (1994), Romania (1996 to 1997), the Czech Republic (1997), Congo (1998), Russia (1999), the United States (1999 to 2009), Canada (1999–2007), Israel (2000), Greece (2010), and Israel (2024). Epizootics of disease in horses occurred in Morocco (1996), Italy (1998), the United States (1999 to 2001), and France (2000), Mexico (2003) and Sardinia (2011). In August 2024 in Warsaw the West Nile virus was identified in bodies of dead birds (Corvidae) while investigating an unusually high number of finds. Outdoor workers (including biological fieldworkers, construction workers, farmers, landscapers, and painters), healthcare personnel, and laboratory personnel who perform necropsies on animals are at risk of contracting WNV. In 2012, the US experienced one of its worst epidemics in which 286 people died, with the state of Texas being hard hit by this virus. Weather Drought has been associated with a higher number of West Nile virus cases in the following year. As drought can decrease fish and other populations that eat mosquito eggs, higher numbers of mosquitoes can result. Higher temperatures are linked to decreased time for replication and increased viral load in birds and mosquitoes. Research A vaccine for horses (ATCvet code: ) based on killed viruses exists; some zoos have given this vaccine to their birds, although its effectiveness is unknown. Dogs and cats show few if any signs of infection. There have been no known cases of direct canine-human or feline-human transmission; although these pets can become infected, it is unlikely they are, in turn, capable of infecting native mosquitoes and thus continuing the disease cycle. AMD3100, which had been proposed as an antiretroviral drug for HIV, has shown promise against West Nile encephalitis. Morpholino antisense oligos conjugated to cell penetrating peptides have been shown to partially protect mice from WNV disease. There have also been attempts to treat infections using ribavirin, intravenous immunoglobulin, or alpha interferon. GenoMed, a U.S. biotech company, has found that blocking angiotensin II can treat the "cytokine storm" of West Nile virus encephalitis as well as other viruses. As of 2019, six vaccines had progressed to human trials but none had been licensed in the United States. Only the two live attenuated vaccines produced strong immunity after a single dose. Dr. Anthony Fauci, former director of the National Institute of Allergy and Infectious Diseases, has urged proactive action, including international collaborations for vaccine and antiviral development, emphasizing that we must not wait for a greater crisis to address this virus. His call for increased public awareness and scientific research followed his own recovery as a victim of the West Nile virus himself, which he most likely contracted in his Washington, D.C.–based backyard from a mosquito bite.	Biology and health sciences	Infectious disease	null
37800	https://en.wikipedia.org/wiki/Dendrochronology	Dendrochronology	Dendrochronology (or tree-ring dating) is the scientific method of dating tree rings (also called growth rings) to the exact year they were formed in a tree. As well as dating them, this can give data for dendroclimatology, the study of climate and atmospheric conditions during different periods in history from the wood of old trees. Dendrochronology derives from the Ancient Greek (), meaning "tree", (), meaning "time", and (), "the study of". Dendrochronology is useful for determining the precise age of samples, especially those that are too recent for radiocarbon dating, which always produces a range rather than an exact date. However, for a precise date of the death of the tree a full sample to the edge is needed, which most trimmed timber will not provide. It also gives data on the timing of events and rates of change in the environment (most prominently climate) and also in wood found in archaeology or works of art and architecture, such as old panel paintings. It is also used as a check in radiocarbon dating to calibrate radiocarbon ages. New growth in trees occurs in a layer of cells near the bark. A tree's growth rate changes in a predictable pattern throughout the year in response to seasonal climate changes, resulting in visible growth rings. Each ring marks a complete cycle of seasons, or one year, in the tree's life. As of 2023, securely dated tree-ring data for Germany and Ireland are available going back 13,910 years. A new method is based on measuring variations in oxygen isotopes in each ring, and this 'isotope dendrochronology' can yield results on samples which are not suitable for traditional dendrochronology due to too few or too similar rings. Some regions have "floating sequences", with gaps which mean that earlier periods can only be approximately dated. As of 2024, only three areas have continuous sequences going back to prehistoric times, the foothills of the Northern Alps, the southwestern United States and the British Isles. Miyake events, which are major spikes in cosmic rays at known dates, are visible in trees rings and can fix the dating of a floating sequence. History The Greek botanist Theophrastus (c. 371 – c. 287 BC) first mentioned that the wood of trees has rings. In his Trattato della Pittura (Treatise on Painting), Leonardo da Vinci (1452–1519) was the first person to mention that trees form rings annually and that their thickness is determined by the conditions under which they grew. In 1737, French investigators Henri-Louis Duhamel du Monceau and Georges-Louis Leclerc de Buffon examined the effect of growing conditions on the shape of tree rings. They found that in 1709, a severe winter produced a distinctly dark tree ring, which served as a reference for subsequent European naturalists. In the U.S., Alexander Catlin Twining (1801–1884) suggested in 1833 that patterns among tree rings could be used to synchronize the dendrochronology of various trees and thereby to reconstruct past climates across entire regions. The English polymath Charles Babbage proposed using dendrochronology to date the remains of trees in peat bogs or even in geological strata (1835, 1838). During the latter half of the nineteenth century, the scientific study of tree rings and the application of dendrochronology began. In 1859, the German-American Jacob Kuechler (1823–1893) used crossdating to examine oaks (Quercus stellata) in order to study the record of climate in western Texas. In 1866, the German botanist, entomologist, and forester Julius Theodor Christian Ratzeburg (1801–1871) observed the effects on tree rings of defoliation caused by insect infestations. By 1882, this observation was already appearing in forestry textbooks. In the 1870s, the Dutch astronomer Jacobus Kapteyn (1851–1922) was using crossdating to reconstruct the climates of the Netherlands and Germany. In 1881, the Swiss-Austrian forester Arthur von Seckendorff-Gudent (1845–1886) was using crossdating. From 1869 to 1901, Robert Hartig (1839–1901), a German professor of forest pathology, wrote a series of papers on the anatomy and ecology of tree rings. In 1892, the Russian physicist (1841–1905) wrote that he had used patterns found in tree rings to predict droughts in 1882 and 1891. During the first half of the twentieth century, the astronomer A. E. Douglass founded the Laboratory of Tree-Ring Research at the University of Arizona. Douglass sought to better understand cycles of sunspot activity and reasoned that changes in solar activity would affect climate patterns on earth, which would subsequently be recorded by tree-ring growth patterns (i.e., sunspots → climate → tree rings). Methods Growth rings Horizontal cross sections cut through the trunk of a tree can reveal growth rings, also referred to as tree rings or annual rings. Growth rings result from new growth in the vascular cambium, a layer of cells near the bark that botanists classify as a lateral meristem; this growth in diameter is known as secondary growth. Visible rings result from the change in growth speed through the seasons of the year; thus, critical for the title method, one ring generally marks the passage of one year in the life of the tree. Removal of the bark of the tree in a particular area may cause deformation of the rings as the plant overgrows the scar. The rings are more visible in trees which have grown in temperate zones, where the seasons differ more markedly. The inner portion of a growth ring forms early in the growing season, when growth is comparatively rapid (hence the wood is less dense) and is known as "early wood" (or "spring wood", or "late-spring wood"); the outer portion is the "late wood" (sometimes termed "summer wood", often being produced in the summer, though sometimes in the autumn) and is denser. Many trees in temperate zones produce one growth-ring each year, with the newest adjacent to the bark. Hence, for the entire period of a tree's life, a year-by-year record or ring pattern builds up that reflects the age of the tree and the climatic conditions in which the tree grew. Adequate moisture and a long growing season result in a wide ring, while a drought year may result in a very narrow one. Direct reading of tree ring chronologies is a complex science, for several reasons. First, contrary to the single-ring-per-year paradigm, alternating poor and favorable conditions, such as mid-summer droughts, can result in several rings forming in a given year. In addition, particular tree species may present "missing rings", and this influences the selection of trees for study of long time-spans. For instance, missing rings are rare in oak and elm trees. Critical to the science, trees from the same region tend to develop the same patterns of ring widths for a given period of chronological study. Researchers can compare and match these patterns ring-for-ring with patterns from trees which have grown at the same time in the same geographical zone (and therefore under similar climatic conditions). When one can match these tree-ring patterns across successive trees in the same locale, in overlapping fashion, chronologies can be built up—both for entire geographical regions and for sub-regions. Moreover, wood from ancient structures with known chronologies can be matched to the tree-ring data (a technique called 'cross-dating'), and the age of the wood can thereby be determined precisely. Dendrochronologists originally carried out cross-dating by visual inspection; more recently, they have harnessed computers to do the task, applying statistical techniques to assess the matching. To eliminate individual variations in tree-ring growth, dendrochronologists take the smoothed average of the tree-ring widths of multiple tree-samples to build up a 'ring history', a process termed replication. A tree-ring history whose beginning- and end-dates are not known is called a 'floating chronology'. It can be anchored by cross-matching a section against another chronology (tree-ring history) whose dates are known. A fully anchored and cross-matched chronology for oak and pine in central Europe extends back 12,460 years, and an oak chronology goes back 7,429 years in Ireland and 6,939 years in England. Comparison of radiocarbon and dendrochronological ages supports the consistency of these two independent dendrochronological sequences. Another fully anchored chronology that extends back 8,500 years exists for the bristlecone pine in the Southwest US (White Mountains of California). Dendrochronological equation The dendrochronological equation defines the law of growth of tree rings. The equation was proposed by Russian biophysicist Alexandr N. Tetearing in his work "Theory of populations" in the form: where ΔL is width of annual ring, t is time (in years), ρ is density of wood, kv is some coefficient, M(t) is function of mass growth of the tree. Ignoring the natural sinusoidal oscillations in tree mass, the formula for the changes in the annual ring width is: where c1, c2, and c4 are some coefficients, a1 and a2 are positive constants. The formula is useful for correct approximation of samples data before data normalization procedure. The typical forms of the function ΔL(t) of annual growth of wood ring are shown in the figures. Sampling and dating Dendrochronology allows specimens of once-living material to be accurately dated to a specific year. Dates are often represented as estimated calendar years B.P., for before present, where "present" refers to 1 January 1950. Timber core samples are sampled and used to measure the width of annual growth rings; by taking samples from different sites within a particular region, researchers can build a comprehensive historical sequence. The techniques of dendrochronology are more consistent in areas where trees grew in marginal conditions such as aridity or semi-aridity where the ring growth is more sensitive to the environment, rather than in humid areas where tree-ring growth is more uniform (complacent). In addition, some genera of trees are more suitable than others for this type of analysis. For instance, the bristlecone pine is exceptionally long-lived and slow growing, and has been used extensively for chronologies; still-living and dead specimens of this species provide tree-ring patterns going back thousands of years, in some regions more than 10,000 years. Currently, the maximum span for fully anchored chronology is a little over 11,000 years B.P. IntCal20 is the 2020 "Radiocarbon Age Calibration Curve", which provides a calibrated carbon 14 dated sequence going back 55,000 years. The most recent part, going back 13,900 years, is based on tree rings. Reference sequences European chronologies derived from wooden structures initially found it difficult to bridge the gap in the fourteenth century when there was a building hiatus, which coincided with the Black Death. However, there do exist unbroken chronologies dating back to prehistoric times, for example the Danish chronology dating back to 352 BC. Given a sample of wood, the variation of the tree-ring growths not only provides a match by year, but can also match location because climate varies from place to place. This makes it possible to determine the source of ships as well as smaller artifacts made from wood, but which were transported long distances, such as panels for paintings and ship timbers. Miyake events Miyake events, such as the ones in 774–775 and 993–994, can provide fixed reference points in an unknown time sequence as they are due to cosmic radiation. As they appear as spikes in carbon 14 in tree rings for that year all round the world, they can be used to date historical events to the year. For example, wooden houses in the Viking site at L'Anse aux Meadows in Newfoundland were dated by finding the layer with the 993 spike, which showed that the wood is from a tree felled in 1021. Researchers at the University of Bern have provided exact dating of a floating sequence in a Neolithic settlement in northern Greece by tying it to a spike in cosmogenic radiocarbon in 5259 BC. Frost rings Frost ring is a term used to designate a layer of deformed, collapsed tracheids and traumatic parenchyma cells in tree ring analysis. They are formed when air temperature falls below freezing during a period of cambial activity. They can be used in dendrochronology to indicate years that are colder than usual. Applications Radiocarbon dating calibration Dates from dendrochronology can be used as a calibration and check of radiocarbon dating. This can be done by checking radiocarbon dates against long master sequences, with Californian bristle-cone pines in Arizona being used to develop this method of calibration as the longevity of the trees (up to c.4900 years) in addition to the use of dead samples meant a long, unbroken tree ring sequence could be developed (dating back to ). Additional studies of European oak trees, such as the master sequence in Germany that dates back to , can also be used to back up and further calibrate radiocarbon dates. Climatology Dendroclimatology is the science of determining past climates from trees primarily from the properties of the annual tree rings. Other properties of the annual rings, such as maximum latewood density (MXD) have been shown to be better proxies than simple ring width. Using tree rings, scientists have estimated many local climates for hundreds to thousands of years previous. Art history Dendrochronology has become important to art historians in the dating of panel paintings. However, unlike analysis of samples from buildings, which are typically sent to a laboratory, wooden supports for paintings usually have to be measured in a museum conservation department, which places limitations on the techniques that can be used. In addition to dating, dendrochronology can also provide information as to the source of the panel. Many Early Netherlandish paintings have turned out to be painted on panels of "Baltic oak" shipped from the Vistula region via ports of the Hanseatic League. Oak panels were used in a number of northern countries such as England, France and Germany. Wooden supports other than oak were rarely used by Netherlandish painters. Since panels of seasoned wood were used, an uncertain number of years has to be allowed for seasoning when estimating dates. Panels were trimmed of the outer rings, and often each panel only uses a small part of the radius of the trunk. Consequently, dating studies usually result in a terminus post quem (earliest possible) date, and a tentative date for the arrival of a seasoned raw panel using assumptions as to these factors. As a result of establishing numerous sequences, it was possible to date 85–90% of the 250 paintings from the fourteenth to seventeenth century analysed between 1971 and 1982; by now a much greater number have been analysed. A portrait of Mary, Queen of Scots in the National Portrait Gallery, London was believed to be an eighteenth-century copy. However, dendrochronology revealed that the wood dated from the second half of the sixteenth century. It is now regarded as an original sixteenth-century painting by an unknown artist. On the other hand, dendrochronology was applied to four paintings depicting the same subject, that of Christ expelling the money-lenders from the Temple. The results showed that the age of the wood was too late for any of them to have been painted by Hieronymus Bosch. While dendrochronology has become an important tool for dating oak panels, it is not effective in dating the poplar panels often used by Italian painters because of the erratic growth rings in poplar. The sixteenth century saw a gradual replacement of wooden panels by canvas as the support for paintings, which means the technique is less often applicable to later paintings. In addition, many panel paintings were transferred onto canvas or other supports during the nineteenth and twentieth centuries. Archaeology The dating of buildings with wooden structures and components is also done by dendrochronology; dendroarchaeology is the term for the application of dendrochronology in archaeology. While archaeologists can date wood and when it was felled, it may be difficult to definitively determine the age of a building or structure in which the wood was used; the wood could have been reused from an older structure, may have been felled and left for many years before use, or could have been used to replace a damaged piece of wood. The dating of building via dendrochronology thus requires knowledge of the history of building technology. Many prehistoric forms of buildings used "posts" that were whole young tree trunks; where the bottom of the post has survived in the ground these can be especially useful for dating. Examples: The Post Track and Sweet Track, ancient timber trackways in the Somerset levels, England, have been dated to 3838 BC and 3807 BC. Navan Fort where in Prehistoric Ireland a large structure was built with more than two hundred posts. The central oak post was felled in 95 BC. The Fairbanks House in Dedham, Massachusetts. While the house had long been claimed to have been built (and being the oldest wood-framed house in North America), core samples of wood taken from a summer beam confirmed the wood was from an oak tree felled in 1637–8, as wood was not seasoned before use in building at that time in New England. An additional sample from another beam yielded a date of 1641, thus confirming the house had been constructed starting in 1638 and finished sometime after 1641 . The burial chamber of Gorm the Old, who died c. 958, was constructed from wood of timbers felled in 958. Veliky Novgorod, where, between the tenth and the fifteenth century, numerous consecutive layers of wooden log pavement have been placed over the accumulating dirt. Measurement platforms, software, and data formats There are many different file formats used to store tree ring width data. Effort for standardisation was made with the development of TRiDaS. Further development led to the database software Tellervo, which is based on the new standard format whilst being able to import lots of different data formats. The desktop application can be attached to measurement devices and works with the database server that is installed separately. Continuous sequence Bard et al write in 2023: "The oldest tree-ring series are known as floating since, while their constituent rings can be counted to create a relative internal chronology, they cannot be dendro-matched with the main Holocene absolute chronology. However, 14C analyses performed at high resolution on overlapped absolute and floating tree-rings series enable one to link them almost absolutely and hence to extend the calibration on annual tree rings until ≈13 900 cal yr BP." Related chronologies Herbchronology is the analysis of annual growth rings (or simply annual rings) in the secondary root xylem of perennial herbaceous plants. Similar seasonal patterns also occur in ice cores and in varves (layers of sediment deposition in a lake, river, or sea bed). The deposition pattern in the core will vary for a frozen-over lake versus an ice-free lake, and with the fineness of the sediment. Sclerochronology is the study of algae deposits. Some columnar cacti also exhibit similar seasonal patterns in the isotopes of carbon and oxygen in their spines (acanthochronology). These are used for dating in a manner similar to dendrochronology, and such techniques are used in combination with dendrochronology, to plug gaps and to extend the range of the seasonal data available to archaeologists and paleoclimatologists. A similar technique is used to estimate the age of fish stocks through the analysis of growth rings in the otolith bones.	Physical sciences	Geology: General	Earth science
37817	https://en.wikipedia.org/wiki/M%C3%B6bius%20strip	Möbius strip	In mathematics, a Möbius strip, Möbius band, or Möbius loop is a surface that can be formed by attaching the ends of a strip of paper together with a half-twist. As a mathematical object, it was discovered by Johann Benedict Listing and August Ferdinand Möbius in 1858, but it had already appeared in Roman mosaics from the third century CE. The Möbius strip is a non-orientable surface, meaning that within it one cannot consistently distinguish clockwise from counterclockwise turns. Every non-orientable surface contains a Möbius strip. As an abstract topological space, the Möbius strip can be embedded into three-dimensional Euclidean space in many different ways: a clockwise half-twist is different from a counterclockwise half-twist, and it can also be embedded with odd numbers of twists greater than one, or with a knotted centerline. Any two embeddings with the same knot for the centerline and the same number and direction of twists are topologically equivalent. All of these embeddings have only one side, but when embedded in other spaces, the Möbius strip may have two sides. It has only a single boundary curve. Several geometric constructions of the Möbius strip provide it with additional structure. It can be swept as a ruled surface by a line segment rotating in a rotating plane, with or without self-crossings. A thin paper strip with its ends joined to form a Möbius strip can bend smoothly as a developable surface or be folded flat; the flattened Möbius strips include the trihexaflexagon. The Sudanese Möbius strip is a minimal surface in a hypersphere, and the Meeks Möbius strip is a self-intersecting minimal surface in ordinary Euclidean space. Both the Sudanese Möbius strip and another self-intersecting Möbius strip, the cross-cap, have a circular boundary. A Möbius strip without its boundary, called an open Möbius strip, can form surfaces of constant curvature. Certain highly symmetric spaces whose points represent lines in the plane have the shape of a Möbius strip. The many applications of Möbius strips include mechanical belts that wear evenly on both sides, dual-track roller coasters whose carriages alternate between the two tracks, and world maps printed so that antipodes appear opposite each other. Möbius strips appear in molecules and devices with novel electrical and electromechanical properties, and have been used to prove impossibility results in social choice theory. In popular culture, Möbius strips appear in artworks by M. C. Escher, Max Bill, and others, and in the design of the recycling symbol. Many architectural concepts have been inspired by the Möbius strip, including the building design for the NASCAR Hall of Fame. Performers including Harry Blackstone Sr. and Thomas Nelson Downs have based stage magic tricks on the properties of the Möbius strip. The canons of J. S. Bach have been analyzed using Möbius strips. Many works of speculative fiction feature Möbius strips; more generally, a plot structure based on the Möbius strip, of events that repeat with a twist, is common in fiction. History The discovery of the Möbius strip as a mathematical object is attributed independently to the German mathematicians Johann Benedict Listing and August Ferdinand Möbius in However, it had been known long before, both as a physical object and in artistic depictions; in particular, it can be seen in several Roman mosaics from the In many cases these merely depict coiled ribbons as boundaries. When the number of coils is odd, these ribbons are Möbius strips, but for an even number of coils they are topologically equivalent to untwisted rings. Therefore, whether the ribbon is a Möbius strip may be coincidental, rather than a deliberate choice. In at least one case, a ribbon with different colors on different sides was drawn with an odd number of coils, forcing its artist to make a clumsy fix at the point where the colors did not Another mosaic from the town of Sentinum (depicted) shows the zodiac, held by the god Aion, as a band with only a single twist. There is no clear evidence that the one-sidedness of this visual representation of celestial time was intentional; it could have been chosen merely as a way to make all of the signs of the zodiac appear on the visible side of the strip. Some other ancient depictions of the ourobouros or of figure-eight-shaped decorations are also alleged to depict Möbius strips, but whether they were intended to depict flat strips of any type is Independently of the mathematical tradition, machinists have long known that mechanical belts wear half as quickly when they form Möbius strips, because they use the entire surface of the belt rather than only the inner surface of an untwisted belt. Additionally, such a belt may be less prone to curling from side to side. An early written description of this technique dates to 1871, which is after the first mathematical publications regarding the Möbius strip. Much earlier, an image of a chain pump in a work of Ismail al-Jazari from 1206 depicts a Möbius strip configuration for its drive Another use of this surface was made by seamstresses in Paris (at an unspecified date): they initiated novices by requiring them to stitch a Möbius strip as a collar onto a Properties The Möbius strip has several curious properties. It is a non-orientable surface: if an asymmetric two-dimensional object slides one time around the strip, it returns to its starting position as its mirror image. In particular, a curved arrow pointing clockwise (↻) would return as an arrow pointing counterclockwise (↺), implying that, within the Möbius strip, it is impossible to consistently define what it means to be clockwise or counterclockwise. It is the simplest non-orientable surface: any other surface is non-orientable if and only if it has a Möbius strip as a Relatedly, when embedded into Euclidean space, the Möbius strip has only one side. A three-dimensional object that slides one time around the surface of the strip is not mirrored, but instead returns to the same point of the strip on what appears locally to be its other side, showing that both positions are really part of a single side. This behavior is different from familiar orientable surfaces in three dimensions such as those modeled by flat sheets of paper, cylindrical drinking straws, or hollow balls, for which one side of the surface is not connected to the other. However, this is a property of its embedding into space rather than an intrinsic property of the Möbius strip itself: there exist other topological spaces in which the Möbius strip can be embedded so that it has two For instance, if the front and back faces of a cube are glued to each other with a left-right mirror reflection, the result is a three-dimensional topological space (the Cartesian product of a Möbius strip with an interval) in which the top and bottom halves of the cube can be separated from each other by a two-sided Möbius In contrast to disks, spheres, and cylinders, for which it is possible to simultaneously embed an uncountable set of disjoint copies into three-dimensional space, only a countable number of Möbius strips can be simultaneously A path along the edge of a Möbius strip, traced until it returns to its starting point on the edge, includes all boundary points of the Möbius strip in a single continuous curve. For a Möbius strip formed by gluing and twisting a rectangle, it has twice the length of the centerline of the strip. In this sense, the Möbius strip is different from an untwisted ring and like a circular disk in having only one A Möbius strip in Euclidean space cannot be moved or stretched into its mirror image; it is a chiral object with right- or Möbius strips with odd numbers of half-twists greater than one, or that are knotted before gluing, are distinct as embedded subsets of three-dimensional space, even though they are all equivalent as two-dimensional topological More precisely, two Möbius strips are equivalently embedded in three-dimensional space when their centerlines determine the same knot and they have the same number of twists as each With an even number of twists, however, one obtains a different topological surface, called the The Möbius strip can be continuously transformed into its centerline, by making it narrower while fixing the points on the centerline. This transformation is an example of a deformation retraction, and its existence means that the Möbius strip has many of the same properties as its centerline, which is topologically a circle. In particular, its fundamental group is the same as the fundamental group of a circle, an infinite cyclic group. Therefore, paths on the Möbius strip that start and end at the same point can be distinguished topologically (up to homotopy) only by the number of times they loop around the strip. Cutting a Möbius strip along the centerline with a pair of scissors yields one long strip with four half-twists in it (relative to an untwisted annulus or cylinder) rather than two separate strips. Two of the half-twists come from the fact that this thinner strip goes two times through the half-twist in the original Möbius strip, and the other two come from the way the two halves of the thinner strip wrap around each other. The result is not a Möbius strip, but instead is topologically equivalent to a cylinder. Cutting this double-twisted strip again along its centerline produces two linked double-twisted strips. If, instead, a Möbius strip is cut lengthwise, a third of the way across its width, it produces two linked strips. One of the two is a central, thinner, Möbius strip, while the other has two These interlinked shapes, formed by lengthwise slices of Möbius strips with varying widths, are sometimes called paradromic The Möbius strip can be cut into six mutually-adjacent regions, showing that maps on the surface of the Möbius strip can sometimes require six colors, in contrast to the four color theorem for the Six colors are always enough. This result is part of the Ringel–Youngs theorem, which states how many colors each topological surface The edges and vertices of these six regions form Tietze's graph, which is a dual graph on this surface for the six-vertex complete graph but cannot be drawn without crossings on a plane. Another family of graphs that can be embedded on the Möbius strip, but not on the plane, are the Möbius ladders, the boundaries of subdivisions of the Möbius strip into rectangles meeting These include the utility graph, a six-vertex complete bipartite graph whose embedding into the Möbius strip shows that, unlike in the plane, the three utilities problem can be solved on a transparent Möbius The Euler characteristic of the Möbius strip is zero, meaning that for any subdivision of the strip by vertices and edges into regions, the numbers , , and of vertices, edges, and regions satisfy . For instance, Tietze's graph has vertices, edges, and regions; Constructions There are many different ways of defining geometric surfaces with the topology of the Möbius strip, yielding realizations with additional geometric properties. Sweeping a line segment One way to embed the Möbius strip in three-dimensional Euclidean space is to sweep it out by a line segment rotating in a plane, which in turn rotates around one of its For the swept surface to meet up with itself after a half-twist, the line segment should rotate around its center at half the angular velocity of the plane's rotation. This can be described as a parametric surface defined by equations for the Cartesian coordinates of its points, for and where one parameter describes the rotation angle of the plane around its central axis and the other parameter describes the position of a point along the rotating line segment. This produces a Möbius strip of width 1, whose center circle has radius 1, lies in the -plane and is centered at The same method can produce Möbius strips with any odd number of half-twists, by rotating the segment more quickly in its plane. The rotating segment sweeps out a circular disk in the plane that it rotates within, and the Möbius strip that it generates forms a slice through the solid torus swept out by this disk. Because of the one-sidedness of this slice, the sliced torus remains A line or line segment swept in a different motion, rotating in a horizontal plane around the origin as it moves up and down, forms Plücker's conoid or cylindroid, an algebraic ruled surface in the form of a self-crossing Möbius It has applications in the design of Polyhedral surfaces and flat foldings A strip of paper can form a flattened Möbius strip in the plane by folding it at angles so that its center line lies along an equilateral triangle, and attaching the ends. The shortest strip for which this is possible consists of three equilateral triangles, folded at the edges where two triangles meet. Its aspect ratiothe ratio of the strip's length to its widthis and the same folding method works for any larger aspect For a strip of nine equilateral triangles, the result is a trihexaflexagon, which can be flexed to reveal different parts of its For strips too short to apply this method directly, one can first "accordion fold" the strip in its wide direction back and forth using an even number of folds. With two folds, for example, a strip would become a folded strip whose cross section is in the shape of an 'N' and would remain an 'N' after a half-twist. The narrower accordion-folded strip can then be folded and joined in the same way that a longer strip The Möbius strip can also be embedded as a polyhedral surface in space or flat-folded in the plane, with only five triangular faces sharing five vertices. In this sense, it is simpler than the cylinder, which requires six triangles and six vertices, even when represented more abstractly as a simplicial complex. A five-triangle Möbius strip can be represented most symmetrically by five of the ten equilateral triangles of a four-dimensional regular simplex. This four-dimensional polyhedral Möbius strip is the only tight Möbius strip, one that is fully four-dimensional and for which all cuts by hyperplanes separate it into two parts that are topologically equivalent to disks or Other polyhedral embeddings of Möbius strips include one with four convex quadrilaterals as faces, another with three non-convex quadrilateral and one using the vertices and center point of a regular octahedron, with a triangular Every abstract triangulation of the projective plane can be embedded into 3D as a polyhedral Möbius strip with a triangular boundary after removing one of its an example is the six-vertex projective plane obtained by adding one vertex to the five-vertex Möbius strip, connected by triangles to each of its boundary However, not every abstract triangulation of the Möbius strip can be represented geometrically, as a polyhedral To be realizable, it is necessary and sufficient that there be no two disjoint non-contractible 3-cycles in the Smoothly embedded rectangles A rectangular Möbius strip, made by attaching the ends of a paper rectangle, can be embedded smoothly into three-dimensional space whenever its aspect ratio is greater than the same ratio as for the flat-folded equilateral-triangle version of the Möbius This flat triangular embedding can lift to a smooth embedding in three dimensions, in which the strip lies flat in three parallel planes between three cylindrical rollers, each tangent to two of the Mathematically, a smoothly embedded sheet of paper can be modeled as a developable surface, that can bend but cannot As its aspect ratio decreases toward , all smooth embeddings seem to approach the same triangular The lengthwise folds of an accordion-folded flat Möbius strip prevent it from forming a three-dimensional embedding in which the layers are separated from each other and bend smoothly without crumpling or stretching away from the Instead, unlike in the flat-folded case, there is a lower limit to the aspect ratio of smooth rectangular Möbius strips. Their aspect ratio cannot be less than even if self-intersections are allowed. Self-intersecting smooth Möbius strips exist for any aspect ratio above this Without self-intersections, the aspect ratio must be at For aspect ratios between this bound it has been an open problem whether smooth embeddings, without self-intersection, In 2023, Richard Schwartz announced a proof that they do not exist, but this result still awaits peer review and publication. If the requirement of smoothness is relaxed to allow continuously differentiable surfaces, the Nash–Kuiper theorem implies that any two opposite edges of any rectangle can be glued to form an embedded Möbius strip, no matter how small the aspect ratio The limiting case, a surface obtained from an infinite strip of the plane between two parallel lines, glued with the opposite orientation to each other, is called the unbounded Möbius strip or the real tautological line bundle. Although it has no smooth closed embedding into three-dimensional space, it can be embedded smoothly as a closed subset of four-dimensional Euclidean The minimum-energy shape of a smooth Möbius strip glued from a rectangle does not have a known analytic description, but can be calculated numerically, and has been the subject of much study in plate theory since the initial work on this subject in 1930 by Michael Sadowsky. It is also possible to find algebraic surfaces that contain rectangular developable Möbius Making the boundary circular The edge, or boundary, of a Möbius strip is topologically equivalent to a circle. In common forms of the Möbius strip, it has a different shape from a circle, but it is unknotted, and therefore the whole strip can be stretched without crossing itself to make the edge perfectly One such example is based on the topology of the Klein bottle, a one-sided surface with no boundary that cannot be embedded into three-dimensional space, but can be immersed (allowing the surface to cross itself in certain restricted ways). A Klein bottle is the surface that results when two Möbius strips are glued together edge-to-edge, andreversing that processa Klein bottle can be sliced along a carefully chosen cut to produce two Möbius For a form of the Klein bottle known as Lawson's Klein bottle, the curve along which it is sliced can be made circular, resulting in Möbius strips with circular Lawson's Klein bottle is a self-crossing minimal surface in the unit hypersphere of 4-dimensional space, the set of points of the form for Half of this Klein bottle, the subset with , gives a Möbius strip embedded in the hypersphere as a minimal surface with a great circle as its This embedding is sometimes called the "Sudanese Möbius strip" after topologists Sue Goodman and Daniel Asimov, who discovered it in the Geometrically Lawson's Klein bottle can be constructed by sweeping a great circle through a great-circular motion in the 3-sphere, and the Sudanese Möbius strip is obtained by sweeping a semicircle instead of a circle, or equivalently by slicing the Klein bottle along a circle that is perpendicular to all of the swept Stereographic projection transforms this shape from a three-dimensional spherical space into three-dimensional Euclidean space, preserving the circularity of its The most symmetric projection is obtained by using a projection point that lies on that great circle that runs through the midpoint of each of the semicircles, but produces an unbounded embedding with the projection point removed from its Instead, leaving the Sudanese Möbius strip unprojected, in the 3-sphere, leaves it with an infinite group of symmetries isomorphic to the orthogonal group the group of symmetries of a The Sudanese Möbius strip extends on all sides of its boundary circle, unavoidably if the surface is to avoid crossing itself. Another form of the Möbius strip, called the cross-cap or crosscap, also has a circular boundary, but otherwise stays on only one side of the plane of this making it more convenient for attaching onto circular holes in other surfaces. In order to do so, it crosses itself. It can be formed by removing a quadrilateral from the top of a hemisphere, orienting the edges of the quadrilateral in alternating directions, and then gluing opposite pairs of these edges consistently with this The two parts of the surface formed by the two glued pairs of edges cross each other with a pinch point like that of a Whitney umbrella at each end of the crossing the same topological structure seen in Plücker's Surfaces of constant curvature The open Möbius strip is the relative interior of a standard Möbius strip, formed by omitting the points on its boundary edge. It may be given a Riemannian geometry of constant positive, negative, or zero Gaussian curvature. The cases of negative and zero curvature form geodesically complete surfaces, which means that all geodesics ("straight lines" on the surface) may be extended indefinitely in either direction. Zero curvature An open strip with zero curvature may be constructed by gluing the opposite sides of a plane strip between two parallel lines, described above as the tautological line The resulting metric makes the open Möbius strip into a (geodesically) complete flat surface (i.e., having zero Gaussian curvature everywhere). This is the unique metric on the Möbius strip, up to uniform scaling, that is both flat and complete. It is the quotient space of a plane by a glide reflection, and (together with the plane, cylinder, torus, and Klein bottle) is one of only five two-dimensional complete Negative curvature The open Möbius strip also admits complete metrics of constant negative curvature. One way to see this is to begin with the upper half plane (Poincaré) model of the hyperbolic plane, a geometry of constant curvature whose lines are represented in the model by semicircles that meet the -axis at right angles. Take the subset of the upper half-plane between any two nested semicircles, and identify the outer semicircle with the left-right reversal of the inner semicircle. The result is topologically a complete and non-compact Möbius strip with constant negative curvature. It is a "nonstandard" complete hyperbolic surface in the sense that it contains a complete hyperbolic half-plane (actually two, on opposite sides of the axis of glide-reflection), and is one of only 13 nonstandard Again, this can be understood as the quotient of the hyperbolic plane by a glide Positive curvature A Möbius strip of constant positive curvature cannot be complete, since it is known that the only complete surfaces of constant positive curvature are the sphere and the projective plane. However, in a sense it is only one point away from being a complete surface, as the open Möbius strip is homeomorphic to the once-punctured projective plane, the surface obtained by removing any one point from the projective The minimal surfaces are described as having constant zero mean curvature instead of constant Gaussian curvature. The Sudanese Möbius strip was constructed as a minimal surface bounded by a great circle in a 3-sphere, but there is also a unique complete (boundaryless) minimal surface immersed in Euclidean space that has the topology of an open Möbius strip. It is called the Meeks Möbius after its 1982 description by William Hamilton Meeks, III. Although globally unstable as a minimal surface, small patches of it, bounded by non-contractible curves within the surface, can form stable embedded Möbius strips as minimal Both the Meeks Möbius strip, and every higher-dimensional minimal surface with the topology of the Möbius strip, can be constructed using solutions to the Björling problem, which defines a minimal surface uniquely from its boundary curve and tangent planes along this Spaces of lines The family of lines in the plane can be given the structure of a smooth space, with each line represented as a point in this space. The resulting space of lines is topologically equivalent to the open Möbius One way to see this is to extend the Euclidean plane to the real projective plane by adding one more line, the line at infinity. By projective duality the space of lines in the projective plane is equivalent to its space of points, the projective plane itself. Removing the line at infinity, to produce the space of Euclidean lines, punctures this space of projective Therefore, the space of Euclidean lines is a punctured projective plane, which is one of the forms of the open Möbius The space of lines in the hyperbolic plane can be parameterized by unordered pairs of distinct points on a circle, the pairs of points at infinity of each line. This space, again, has the topology of an open Möbius These spaces of lines are highly symmetric. The symmetries of Euclidean lines include the affine transformations, and the symmetries of hyperbolic lines include the The affine transformations and Möbius transformations both form Lie groups, topological spaces having a compatible algebraic structure describing the composition of Because every line in the plane is symmetric to every other line, the open Möbius strip is a homogeneous space, a space with symmetries that take every point to every other point. Homogeneous spaces of Lie groups are called solvmanifolds, and the Möbius strip can be used as a counterexample, showing that not every solvmanifold is a nilmanifold, and that not every solvmanifold can be factored into a direct product of a compact solvmanifold These symmetries also provide another way to construct the Möbius strip itself, as a group model of these Lie groups. A group model consists of a Lie group and a stabilizer subgroup of its action; contracting the cosets of the subgroup to points produces a space with the same topology as the underlying homogenous space. In the case of the symmetries of Euclidean lines, the stabilizer of the consists of all symmetries that take the axis to itself. Each line corresponds to a coset, the set of symmetries that map to the Therefore, the quotient space, a space that has one point per coset and inherits its topology from the space of symmetries, is the same as the space of lines, and is again an open Möbius Applications Beyond the already-discussed applications of Möbius strips to the design of mechanical belts that wear evenly on their entire surface, and of the Plücker conoid to the design of gears, other applications of Möbius strips include: Graphene ribbons twisted to form Möbius strips with new electronic characteristics including helical magnetism Möbius aromaticity, a property of organic chemicals whose molecular structure forms a cycle, with molecular orbitals aligned along the cycle in the pattern of a Möbius strip The Möbius resistor, a strip of conductive material covering the single side of a dielectric Möbius strip, in a way that cancels its own self-inductance Resonators with a compact design and a resonant frequency that is half that of identically constructed linear coils Polarization patterns in light emerging from a q-plate A proof of the impossibility of continuous, anonymous, and unanimous two-party aggregation rules in social choice theory Möbius loop roller coasters, a form of dual-tracked roller coaster in which the two tracks spiral around each other an odd number of times, so that the carriages return to the other track than the one they started on World maps projected onto a Möbius strip with the convenient properties that there are no east–west boundaries, and that the antipode of any point on the map can be found on the other printed side of the surface at the same point of the Möbius strip Scientists have also studied the energetics of soap films shaped as Möbius strips, the chemical synthesis of molecules with a Möbius strip shape, and the formation of larger nanoscale Möbius strips using DNA origami. In popular culture Two-dimensional artworks featuring the Möbius strip include an untitled 1947 painting by Corrado Cagli (memorialized in a poem by Charles Olson), and two prints by M. C. Escher: Möbius Band I (1961), depicting three folded flatfish biting each others' tails; and Möbius Band II (1963), depicting ants crawling around a lemniscate-shaped Möbius strip. It is also a popular subject of mathematical sculpture, including works by Max Bill (Endless Ribbon, 1953), José de Rivera (Infinity, 1967), and Sebastián. A trefoil-knotted Möbius strip was used in John Robinsons Immortality (1982). Charles O. Perry's Continuum (1976) is one of several pieces by Perry exploring variations of the Möbius strip. Because of their easily recognized form, Möbius strips are a common element of graphic design. The familiar three-arrow logo for recycling, designed in 1970, is based on the smooth triangular form of the Möbius as was the logo for the environmentally-themed Expo '74. Some variations of the recycling symbol use a different embedding with three half-twists instead of and the original version of the Google Drive logo used a flat-folded three-twist Möbius strip, as have other similar designs. The Brazilian Instituto Nacional de Matemática Pura e Aplicada (IMPA) uses a stylized smooth Möbius strip as their logo, and has a matching large sculpture of a Möbius strip on display in their building. The Möbius strip has also featured in the artwork for postage stamps from countries including Brazil, Belgium, the Netherlands, and Möbius strips have been a frequent inspiration for the architectural design of buildings and bridges. However, many of these are projects or conceptual designs rather than constructed objects, or stretch their interpretation of the Möbius strip beyond its recognizability as a mathematical form or a functional part of the architecture. An example is the National Library of Kazakhstan, for which a building was planned in the shape of a thickened Möbius strip but refinished with a different design after the original architects pulled out of the project. One notable building incorporating a Möbius strip is the NASCAR Hall of Fame, which is surrounded by a large twisted ribbon of stainless steel acting as a façade and canopy, and evoking the curved shapes of racing tracks. On a smaller scale, Moebius Chair (2006) by Pedro Reyes is a courting bench whose base and sides have the form of a Möbius strip. As a form of mathematics and fiber arts, scarves have been knit into Möbius strips since the work of Elizabeth Zimmermann in the early 1980s. In food styling, Möbius strips have been used for slicing bagels, making loops out of bacon, and creating new shapes for pasta. Although mathematically the Möbius strip and the fourth dimension are both purely spatial concepts, they have often been invoked in speculative fiction as the basis for a time loop into which unwary victims may become trapped. Examples of this trope include Martin Gardners "No-Sided Professor" (1946), Armin Joseph Deutschs "A Subway Named Mobius" (1950) and the film Moebius (1996) based on it. An entire world shaped like a Möbius strip is the setting of Arthur C. Clarke's "The Wall of Darkness" (1946), while conventional Möbius strips are used as clever inventions in multiple stories of William Hazlett Upson from the 1940s. Other works of fiction have been analyzed as having a Möbius strip–like structure, in which elements of the plot repeat with a twist; these include Marcel Proust In Search of Lost Time (1913–1927), Luigi Pirandello Six Characters in Search of an Author (1921), Frank Capras It's a Wonderful Life (1946), John Barth Lost in the Funhouse (1968), Samuel R. Delanys Dhalgren (1975) and the film Donnie Darko (2001). One of the musical canons by J. S. Bach, the fifth of 14 canons (BWV 1087) discovered in 1974 in Bach's copy of the Goldberg Variations, features a glide-reflect symmetry in which each voice in the canon repeats, with inverted notes, the same motif from two measures earlier. Because of this symmetry, this canon can be thought of as having its score written on a Möbius strip. In music theory, tones that differ by an octave are generally considered to be equivalent notes, and the space of possible notes forms a circle, the chromatic circle. Because the Möbius strip is the configuration space of two unordered points on a circle, the space of all two-note chords takes the shape of a Möbius strip. This conception, and generalizations to more points, is a significant application of orbifolds to music theory. Modern musical groups taking their name from the Möbius strip include American electronic rock trio Mobius Band and Norwegian progressive rock band Ring Van Möbius. Möbius strips and their properties have been used in the design of stage magic. One such trick, known as the Afghan bands, uses the fact that the Möbius strip remains in one piece as a single strip when cut lengthwise. It originated in the 1880s, and was very popular in the first half of the twentieth century. Many versions of this trick exist and have been performed by famous illusionists such as Harry Blackstone Sr. and Thomas Nelson Downs.	Mathematics	Three-dimensional space	null
37864	https://en.wikipedia.org/wiki/Nyquist%E2%80%93Shannon%20sampling%20theorem	Nyquist–Shannon sampling theorem	The Nyquist–Shannon sampling theorem is an essential principle for digital signal processing linking the frequency range of a signal and the sample rate required to avoid a type of distortion called aliasing. The theorem states that the sample rate must be at least twice the bandwidth of the signal to avoid aliasing. In practice, it is used to select band-limiting filters to keep aliasing below an acceptable amount when an analog signal is sampled or when sample rates are changed within a digital signal processing function. The Nyquist–Shannon sampling theorem is a theorem in the field of signal processing which serves as a fundamental bridge between continuous-time signals and discrete-time signals. It establishes a sufficient condition for a sample rate that permits a discrete sequence of samples to capture all the information from a continuous-time signal of finite bandwidth. Strictly speaking, the theorem only applies to a class of mathematical functions having a Fourier transform that is zero outside of a finite region of frequencies. Intuitively we expect that when one reduces a continuous function to a discrete sequence and interpolates back to a continuous function, the fidelity of the result depends on the density (or sample rate) of the original samples. The sampling theorem introduces the concept of a sample rate that is sufficient for perfect fidelity for the class of functions that are band-limited to a given bandwidth, such that no actual information is lost in the sampling process. It expresses the sufficient sample rate in terms of the bandwidth for the class of functions. The theorem also leads to a formula for perfectly reconstructing the original continuous-time function from the samples. Perfect reconstruction may still be possible when the sample-rate criterion is not satisfied, provided other constraints on the signal are known (see below and compressed sensing). In some cases (when the sample-rate criterion is not satisfied), utilizing additional constraints allows for approximate reconstructions. The fidelity of these reconstructions can be verified and quantified utilizing Bochner's theorem. The name Nyquist–Shannon sampling theorem honours Harry Nyquist and Claude Shannon, but the theorem was also previously discovered by E. T. Whittaker (published in 1915), and Shannon cited Whittaker's paper in his work. The theorem is thus also known by the names Whittaker–Shannon sampling theorem, Whittaker–Shannon, and Whittaker–Nyquist–Shannon, and may also be referred to as the cardinal theorem of interpolation. Introduction Sampling is a process of converting a signal (for example, a function of continuous time or space) into a sequence of values (a function of discrete time or space). Shannon's version of the theorem states: A sufficient sample-rate is therefore anything larger than samples per second. Equivalently, for a given sample rate , perfect reconstruction is guaranteed possible for a bandlimit . When the bandlimit is too high (or there is no bandlimit), the reconstruction exhibits imperfections known as aliasing. Modern statements of the theorem are sometimes careful to explicitly state that must contain no sinusoidal component at exactly frequency or that must be strictly less than one half the sample rate. The threshold is called the Nyquist rate and is an attribute of the continuous-time input to be sampled. The sample rate must exceed the Nyquist rate for the samples to suffice to represent The threshold is called the Nyquist frequency and is an attribute of the sampling equipment. All meaningful frequency components of the properly sampled exist below the Nyquist frequency. The condition described by these inequalities is called the Nyquist criterion, or sometimes the Raabe condition. The theorem is also applicable to functions of other domains, such as space, in the case of a digitized image. The only change, in the case of other domains, is the units of measure attributed to and The symbol is customarily used to represent the interval between samples and is called the sample period or sampling interval. The samples of function are commonly denoted by (alternatively in older signal processing literature), for all integer values of The multiplier is a result of the transition from continuous time to discrete time (see Discrete-time Fourier transform#Relation to Fourier Transform), and it preserves the energy of the signal as varies. A mathematically ideal way to interpolate the sequence involves the use of sinc functions. Each sample in the sequence is replaced by a sinc function, centered on the time axis at the original location of the sample with the amplitude of the sinc function scaled to the sample value, Subsequently, the sinc functions are summed into a continuous function. A mathematically equivalent method uses the Dirac comb and proceeds by convolving one sinc function with a series of Dirac delta pulses, weighted by the sample values. Neither method is numerically practical. Instead, some type of approximation of the sinc functions, finite in length, is used. The imperfections attributable to the approximation are known as interpolation error. Practical digital-to-analog converters produce neither scaled and delayed sinc functions, nor ideal Dirac pulses. Instead they produce a piecewise-constant sequence of scaled and delayed rectangular pulses (the zero-order hold), usually followed by a lowpass filter (called an "anti-imaging filter") to remove spurious high-frequency replicas (images) of the original baseband signal. Aliasing When is a function with a Fourier transform : Then the samples, of are sufficient to create a periodic summation of (see Discrete-time Fourier transform#Relation to Fourier Transform): which is a periodic function and its equivalent representation as a Fourier series, whose coefficients are . This function is also known as the discrete-time Fourier transform (DTFT) of the sample sequence. As depicted, copies of are shifted by multiples of the sampling rate and combined by addition. For a band-limited function and sufficiently large it is possible for the copies to remain distinct from each other. But if the Nyquist criterion is not satisfied, adjacent copies overlap, and it is not possible in general to discern an unambiguous Any frequency component above is indistinguishable from a lower-frequency component, called an alias, associated with one of the copies. In such cases, the customary interpolation techniques produce the alias, rather than the original component. When the sample-rate is pre-determined by other considerations (such as an industry standard), is usually filtered to reduce its high frequencies to acceptable levels before it is sampled. The type of filter required is a lowpass filter, and in this application it is called an anti-aliasing filter. Derivation as a special case of Poisson summation When there is no overlap of the copies (also known as "images") of , the term of can be recovered by the product: where: The sampling theorem is proved since uniquely determines . All that remains is to derive the formula for reconstruction. need not be precisely defined in the region because is zero in that region. However, the worst case is when the Nyquist frequency. A function that is sufficient for that and all less severe cases is: where is the rectangular function. Therefore: (from , above). The inverse transform of both sides produces the Whittaker–Shannon interpolation formula: which shows how the samples, , can be combined to reconstruct . Larger-than-necessary values of (smaller values of ), called oversampling, have no effect on the outcome of the reconstruction and have the benefit of leaving room for a transition band in which is free to take intermediate values. Undersampling, which causes aliasing, is not in general a reversible operation. Theoretically, the interpolation formula can be implemented as a low-pass filter, whose impulse response is and whose input is which is a Dirac comb function modulated by the signal samples. Practical digital-to-analog converters (DAC) implement an approximation like the zero-order hold. In that case, oversampling can reduce the approximation error. Shannon's original proof Poisson shows that the Fourier series in produces the periodic summation of , regardless of and . Shannon, however, only derives the series coefficients for the case . Virtually quoting Shannon's original paper: Let be the spectrum of Then because is assumed to be zero outside the band If we let where is any positive or negative integer, we obtain: On the left are values of at the sampling points. The integral on the right will be recognized as essentially the coefficient in a Fourier-series expansion of the function taking the interval to as a fundamental period. This means that the values of the samples determine the Fourier coefficients in the series expansion of Thus they determine since is zero for frequencies greater than and for lower frequencies is determined if its Fourier coefficients are determined. But determines the original function completely, since a function is determined if its spectrum is known. Therefore the original samples determine the function completely. Shannon's proof of the theorem is complete at that point, but he goes on to discuss reconstruction via sinc functions, what we now call the Whittaker–Shannon interpolation formula as discussed above. He does not derive or prove the properties of the sinc function, as the Fourier pair relationship between the rect (the rectangular function) and sinc functions was well known by that time. As in the other proof, the existence of the Fourier transform of the original signal is assumed, so the proof does not say whether the sampling theorem extends to bandlimited stationary random processes.	Technology	Signal processing	null
37872	https://en.wikipedia.org/wiki/Beta%20particle	Beta particle	A beta particle, also called beta ray or beta radiation (symbol β), is a high-energy, high-speed electron or positron emitted by the radioactive decay of an atomic nucleus, known as beta decay. There are two forms of beta decay, β− decay and β+ decay, which produce electrons and positrons, respectively. Beta particles with an energy of 0.5 MeV have a range of about one metre in the air; the distance is dependent on the particle's energy and the air's density and composition. Beta particles are a type of ionizing radiation, and for radiation protection purposes, they are regarded as being more ionising than gamma rays, but less ionising than alpha particles. The higher the ionising effect, the greater the damage to living tissue, but also the lower the penetrating power of the radiation through matter. Beta decay modes β− decay (electron emission) An unstable atomic nucleus with an excess of neutrons may undergo β− decay, where a neutron is converted into a proton, an electron, and an electron antineutrino (the antiparticle of the neutrino): → + + This process is mediated by the weak interaction. The neutron turns into a proton through the emission of a virtual W− boson. At the quark level, W− emission turns a down quark into an up quark, turning a neutron (one up quark and two down quarks) into a proton (two up quarks and one down quark). The virtual W− boson then decays into an electron and an antineutrino. β− decay commonly occurs among the neutron-rich fission byproducts produced in nuclear reactors. Free neutrons also decay via this process. Both of these processes contribute to the copious quantities of beta rays and electron antineutrinos produced by fission-reactor fuel rods. β+ decay (positron emission) Unstable atomic nuclei with an excess of protons may undergo β+ decay, also called positron decay, where a proton is converted into a neutron, a positron, and an electron neutrino: → + + Beta-plus decay can only happen inside nuclei when the absolute value of the binding energy of the daughter nucleus is greater than that of the parent nucleus, i.e., the daughter nucleus is a lower-energy state. Beta decay schemes The accompanying decay scheme diagram shows the beta decay of caesium-137. 137Cs is noted for a characteristic gamma peak at 661 keV, but this is actually emitted by the daughter radionuclide 137mBa. The diagram shows the type and energy of the emitted radiation, its relative abundance, and the daughter nuclides after decay. Phosphorus-32 is a beta emitter widely used in medicine. It has a short half-life of 14.29 days and decays into sulfur-32 by beta decay as shown in this nuclear equation: {\| border="0" \|- style="height:2em;" \| \|\|→ \|\| \|\|+ \|\| \|\|+ \|\| \|} 1.709 MeV of energy is released during the decay. The kinetic energy of the electron varies with an average of approximately 0.5 MeV and the remainder of the energy is carried by the nearly undetectable electron antineutrino. In comparison to other beta radiation-emitting nuclides, the electron is moderately energetic. It is blocked by around 1 m of air or 5 mm of acrylic glass. Interaction with other matter Of the three common types of radiation given off by radioactive materials, alpha, beta and gamma, beta has the medium penetrating power and the medium ionising power. Although the beta particles given off by different radioactive materials vary in energy, most beta particles can be stopped by a few millimeters of aluminium. However, this does not mean that beta-emitting isotopes can be completely shielded by such thin shields: as they decelerate in matter, beta electrons emit secondary gamma rays, which are more penetrating than betas per se. Shielding composed of materials with lower atomic weight generates gammas with lower energy, making such shields somewhat more effective per unit mass than ones made of larger atoms such as lead. Being composed of charged particles, beta radiation is more strongly ionizing than gamma radiation. When passing through matter, a beta particle is decelerated by electromagnetic interactions and may give off bremsstrahlung X-rays. In water, beta radiation from many nuclear fission products typically exceeds the speed of light in that material (which is about 75% that of light in vacuum), and thus generates blue Cherenkov radiation when it passes through water. The intense beta radiation from the fuel rods of swimming pool reactors can thus be visualized through the transparent water that covers and shields the reactor (see illustration at right). Detection and measurement The ionizing or excitation effects of beta particles on matter are the fundamental processes by which radiometric detection instruments detect and measure beta radiation. The ionization of gas is used in ion chambers and Geiger–Müller counters, and the excitation of scintillators is used in scintillation counters. The following table shows radiation quantities in SI and non-SI units: The gray (Gy) is the SI unit of absorbed dose, which is the amount of radiation energy deposited in the irradiated material. For beta radiation this is numerically equal to the equivalent dose measured by the sievert, which indicates the stochastic biological effect of low levels of radiation on human tissue. The radiation weighting conversion factor from absorbed dose to equivalent dose is 1 for beta, whereas alpha particles have a factor of 20, reflecting their greater ionising effect on tissue. The rad is the deprecated CGS unit for absorbed dose and the rem is the deprecated CGS unit of equivalent dose, used mainly in the USA. Beta spectroscopy The energy contained within individual beta particles is measured via beta spectrometry; the study of the obtained distribution of energies as a spectrum is beta spectroscopy. Determination of this energy is done by measuring the amount of deflection of the electron's path under a magnetic field. Applications Beta particles can be used to treat health conditions such as eye and bone cancer and are also used as tracers. Strontium-90 is the material most commonly used to produce beta particles. Beta particles are also used in quality control to test the thickness of an item, such as paper, coming through a system of rollers. Some of the beta radiation is absorbed while passing through the product. If the product is made too thick or thin, a correspondingly different amount of radiation will be absorbed. A computer program monitoring the quality of the manufactured paper will then move the rollers to change the thickness of the final product. An illumination device called a betalight contains tritium and a phosphor. As tritium decays, it emits beta particles; these strike the phosphor, causing the phosphor to give off photons, much like the cathode-ray tube in a television. The illumination requires no external power, and will continue as long as the tritium exists (and the phosphors do not themselves chemically change); the amount of light produced will drop to half its original value in 12.32 years, the half-life of tritium. Beta-plus (or positron) decay of a radioactive tracer isotope is the source of the positrons used in positron emission tomography (PET scan). History Henri Becquerel, while experimenting with fluorescence, accidentally found out that uranium exposed a photographic plate, wrapped with black paper, with some unknown radiation that could not be turned off like X-rays. Ernest Rutherford continued these experiments and discovered two different kinds of radiation: alpha particles that did not show up on the Becquerel plates because they were easily absorbed by the black wrapping paper beta particles which are 100 times more penetrating than alpha particles. He published his results in 1899. In 1900, Becquerel measured the mass-to-charge ratio () for beta particles by the method of J. J. Thomson used to study cathode rays and identify the electron. He found that for a beta particle is the same as for Thomson's electron, and therefore suggested that the beta particle is in fact an electron. Health Beta particles are moderately penetrating in living tissue, and can cause spontaneous mutation in DNA. Beta sources can be used in radiation therapy to kill cancer cells.	Physical sciences	Nuclear physics	null
37882	https://en.wikipedia.org/wiki/Crocodile	Crocodile	Crocodiles (family Crocodylidae) or true crocodiles are large, semiaquatic reptiles that live throughout the tropics in Africa, Asia, the Americas and Australia. The term “crocodile” is sometimes used more loosely to include all extant members of the order Crocodilia, which includes the alligators and caimans (both members of the family Alligatoridae), the gharial and false gharial (both members of the family Gavialidae) as well as other, extinct, taxa. Although crocodiles, alligators, and the gharial are similar in appearance, they belong to separate biological families. The gharial, with its narrow snout, is easier to distinguish, while morphological differences are more difficult to spot in crocodiles and alligators. The most obvious external differences are visible in the head, with crocodiles having narrower and longer heads, with a more V-shaped than a U-shaped snout compared to alligators and caimans. Another obvious trait is that the upper and lower jaws of the crocodiles are the same width, and the teeth in the lower jaw fall along the edge or outside the upper jaw when the mouth is closed; therefore, all teeth are visible, unlike an alligator, which possesses in the upper jaw small depressions into which the lower teeth fit. Also, when the crocodile's mouth is closed, the large fourth tooth in the lower jaw fits into a constriction in the upper jaw. For hard-to-distinguish specimens, the protruding tooth is the most reliable feature to define the species' family. Crocodiles have more webbing on the toes of the hind feet and can better tolerate saltwater due to specialized salt glands for filtering out salt, which are present, but non-functioning, in alligators. Another trait that separates crocodiles from other crocodilians is their much higher levels of aggression. Crocodile size, morphology, behaviour and ecology differ somewhat among species. However, they have many similarities in these areas as well. All crocodiles are semiaquatic and tend to congregate in freshwater habitats such as rivers, lakes, wetlands and sometimes in brackish water and saltwater. They are carnivorous animals, feeding mostly on vertebrates such as fish, reptiles, birds and mammals, and sometimes on invertebrates such as molluscs and crustaceans, depending on species and age. All crocodiles are tropical species that, unlike alligators, are very sensitive to cold. They separated from other crocodilians during the Eocene epoch, about 55 million years ago. Many species are at the risk of extinction, some being classified as critically endangered. Etymology The word crocodile comes , used in the phrase , . There are several variant Greek forms of the word attested, including the later form () found cited in many English reference works. In the Koine Greek of Roman times, and would have been pronounced identically, and either or both may be the source of the Latinized form used by the ancient Romans. It has been suggested, but it is not certain that the word or is a compound of (), and (), although is only attested as a colloquial term for . It is ascribed to Herodotus, and supposedly describes the basking habits of the Egyptian crocodile. The form is attested in Medieval Latin. It is not clear whether this is a medieval corruption or derives from alternative Greco-Latin forms (late Greek and are attested). A (further) corrupted form is found in Old French and was borrowed into Middle English as . The Modern English form crocodile was adapted directly from the Classical Latin in the 16th century, replacing the earlier form. The use of -y- in the scientific name Crocodylus (and forms derived from it) is a corruption introduced by Laurenti (1768). Taxonomy and phylogeny Crocodylidae was named as a family by Georges Cuvier in 1807. It belongs to the larger superfamily Crocodyloidea, which also includes additional extinct crocodile relatives. These all belong to the order Crocodilia, which also includes alligators and gharials. Crocodylidae is cladistically defined as a crown group composed of the last common ancestor of the Nile crocodile (Crocodylus niloticus), the Dwarf crocodile (Osteolaemus tetraspis), and all of its descendants. It contains two subfamilies: Crocodylinae and Osteolaeminae. Crocodylinae contains 13-14 living species, as well as 6 extinct species. Osteolaeminae was named by Christopher Brochu in 2003 as a subfamily of Crocodylidae separate from Crocodylinae, and contains the two extant genera Osteolaemus and Mecistops, along with several extinct genera. The number of extant species within Osteolaeminae is currently in question. Subfamily Crocodylinae Genus Crocodylus Crocodylus acutus, American crocodile Crocodylus halli, Hall's New Guinea crocodile found South of the New Guinea Highlands Crocodylus intermedius, Orinoco crocodile Crocodylus johnsoni, freshwater crocodile, or Johnstone's crocodile Crocodylus mindorensis, Philippine crocodile Crocodylus moreletii, Morelet's crocodile or Mexican crocodile Crocodylus niloticus, Nile crocodile or African crocodile (the subspecies found in Madagascar is sometimes called the black crocodile) Crocodylus novaeguineae, New Guinea crocodile found North of the New Guinea Highlands Crocodylus palustris, mugger, marsh or Indian crocodile Crocodylus porosus, saltwater crocodile or estuarine crocodile Crocodylus raninus, the Borneo crocodile, is currently considered to be a synonym of Crocodylus porosus; whether or not it is a distinct species remains unclear. Crocodylus rhombifer, Cuban crocodile Crocodylus siamensis, Siamese crocodile (may be extinct in the wild) Crocodylus suchus, West African crocodile, desert or sacred crocodile Crocodylus anthropophagus Crocodylus checchiai Crocodylus falconensis Crocodylus palaeindicus Crocodylus thorbjarnarsoni Genus Voay Voay robustus (formerly Crocodylus robustus) Subfamily Osteolaeminae Genus Osteolaemus Osteolaemus tetraspis, dwarf crocodile (There has been controversy as to whether or not this is actually two species; recent (2010) DNA analysis indicate three distinct species: O. tetraspis, O. osborni and a third, currently unnamed.) Genus Mecistops Mecistops cataphractus West African slender-snouted crocodile Mecistops leptorhynchus Central African slender-snouted crocodile Genus Brochuchus Brochuchus pigotti (formerly Crocodylus pigotti) Brochuchus parvidens Genus Euthecodon Euthecodon nitriae Euthecodon brumpti Euthecodon arambourgi Genus Rimasuchus Rimasuchus lloydi (formerly Crocodylus lloydi) Phylogeny Recent molecular studies using DNA sequencing have shown crocodiles to be more closely related to the gavialids rather than to alligators, contrary to prior theories based on morphological studies alone. Below is a cladogram showing the relationships of the major extant crocodile groups based on molecular studies, excluding separate extinct taxa: Below is a more detailed cladogram of Crocodylidae, based on a 2021 study using paleogenomics that extracted DNA from the extinct Voay. Alternatively, some morphological studies have recovered Mecistops as a basal member of Crocodylinae, more closely related to Crocodylus than to Osteolaemus and the other members of Osteolaeminae, as shown in the cladogram below. Species A total of 18 extant species have been recognized. Further genetic study is needed for the confirmation of proposed species under the genus Osteolaemus. Characteristics A crocodile's physical traits allow it to be a successful predator. Its external morphology is a sign of its aquatic and predatory lifestyle. Its streamlined body enables it to swim swiftly; it also tucks its feet to the side while swimming, making it faster by decreasing water resistance. Crocodiles have webbed feet which, though not used to propel them through the water, allow them to make fast turns and sudden moves in the water or initiate swimming. Webbed feet are an advantage in shallow water, where the animals sometimes move around by walking. Crocodiles have a palatal flap, a rigid tissue at the back of the mouth that blocks the entry of water. The palate has a special path from the nostril to the glottis that bypasses the mouth. The nostrils are closed during submergence. Like other archosaurs, crocodilians are diapsid, although their post-temporal fenestrae are reduced. The walls of the braincase are bony but lack supratemporal and postfrontal bones. Their tongues are not free, but held in place by a membrane that limits movement; as a result, crocodiles are unable to stick out their tongues. Crocodiles have smooth skin on their bellies and sides, while their dorsal surfaces are armoured with large osteoderms. The armoured skin has scales and is thick and rugged, providing some protection. They are still able to absorb heat through this armour, as a network of small capillaries allows blood through the scales to absorb heat. The osteoderms are highly vascularised and aid in calcium balance, both to neutralize acids while the animal cannot breathe underwater and to provide calcium for eggshell formation. Crocodilian tegument have pores believed to be sensory in function, analogous to the lateral line in fishes. They are particularly seen on their upper and lower jaws. Another possibility is that they are secretory, as they produce an oily substance which appears to flush mud off. Size Size greatly varies among species, from the dwarf crocodile to the saltwater crocodile. Species of the dwarf crocodile Osteolaemus grow to an adult size of just , whereas the saltwater crocodile can grow to sizes over and weigh over . Several other large species can reach over long and weigh over . Crocodilians show pronounced sexual dimorphism, with males growing much larger and more rapidly than females. Despite their large adult sizes, crocodiles start their lives at around long. The largest species of crocodile is the saltwater crocodile, found in eastern India, northern Australia, throughout South-east Asia, and in the surrounding waters. The brain volume of two adult crocodiles was 5.6 cm3 for a spectacled caiman and 8.5 cm3 for a larger Nile crocodile. The largest crocodile ever held in captivity is a saltwater–Siamese hybrid named Yai (, meaning big; born 10 June 1972) at the Samutprakarn Crocodile Farm and Zoo, Thailand. This animal measures in length and weighs . The longest crocodile captured alive was Lolong, a saltwater crocodile which was measured at and weighed at by a National Geographic team in Agusan del Sur Province, Philippines. Teeth Crocodiles are polyphyodonts; they are able to replace each of their 80 teeth up to 50 times in their 35- to 75-year lifespan. Next to each full-grown tooth, there is a small replacement tooth and an odontogenic stem cell in the dental lamina in standby that can be activated if required. Biology and behaviour Crocodilians are more closely related to birds and dinosaurs than to most animals classified as reptiles, the three families being included in the group Archosauria ('ruling reptiles'). Despite their prehistoric look, crocodiles are among the more biologically complex reptiles. Unlike other reptiles, a crocodile has a cerebral cortex and a four-chambered heart. Crocodilians also have the functional equivalent of a diaphragm by incorporating muscles used for aquatic locomotion into respiration. Salt glands are present in the tongues of crocodiles and they have a pore opening on the surface of the tongue, a trait that separates them from alligators. Salt glands are dysfunctional in Alligatoridae. Their function appears to be similar to that of salt glands in marine turtles. Crocodiles do not have sweat glands and release heat through their mouths. They often sleep with their mouths open and may pant like a dog. Four species of freshwater crocodile climb trees to bask in areas lacking a shoreline. Senses Crocodiles have acute senses, an evolutionary advantage that makes them successful predators. The eyes, ears and nostrils are located on top of the head, allowing the crocodile to lie low in the water, almost totally submerged and hidden from prey. Vision Crocodiles have very good night vision, and are mostly nocturnal hunters. They use the disadvantage of most prey animals' poor nocturnal vision to their advantage. The light receptors in crocodilians' eyes include cones and numerous rods, so it is assumed all crocodilians can see colours. Crocodiles have vertical-slit shaped pupils, similar to those of domestic cats. One explanation for the evolution of slit pupils is that they exclude light more effectively than a circular pupil, helping to protect the eyes during daylight. On the rear wall of the eye is a tapetum lucidum, which reflects incoming light back onto the retina, thus utilizing the small amount of light available at night to best advantage. In addition to the protection of the upper and lower eyelids, crocodiles have a nictitating membrane (sometimes called a "third eye-lid") that can be drawn over the eye from the inner corner while the lids are open. The eyeball surface is thus protected under the water while a certain degree of vision is still possible. Olfaction Crocodilian sense of smell is also very well developed, aiding them to detect prey or animal carcasses that are either on land or in water, from far away. It is possible that crocodiles use olfaction in the egg prior to hatching. Chemoreception in crocodiles is especially interesting because they hunt in both terrestrial and aquatic surroundings. Crocodiles have only one olfactory chamber and the vomeronasal organ is absent in the adults indicating all olfactory perception is limited to the olfactory system. Behavioural and olfactometer experiments indicate that crocodiles detect both air-borne and water-soluble chemicals and use their olfactory system for hunting. When above water, crocodiles enhance their ability to detect volatile odorants by gular pumping, a rhythmic movement of the floor of the pharynx. Crocodiles close their nostrils when submerged, so olfaction underwater is unlikely. Underwater food detection is presumably gustatory and tactile. Hearing Crocodiles can hear well; their tympanic membranes are concealed by flat flaps that may be raised or lowered by muscles. Touch The touch sensors, concentrated in crocodile skin, can be thicker than those in human fingerprints. Crocodiles can feel the touch on their skin. Cranial: The upper and lower jaws are covered with sensory pits, visible as small, black speckles on the skin, the crocodilian version of the lateral line organs seen in fish and many amphibians, though arising from a completely different origin. These pigmented nodules encase bundles of nerve fibers innervated beneath by branches of the trigeminal nerve. They respond to the slightest disturbance in surface water, detecting vibrations and small pressure changes as small as a single drop. This makes it possible for crocodiles to detect prey, danger and intruders, even in total darkness. These sense organs are known as domed pressure receptors (DPRs). Post-Cranial: While alligators and caimans have DPRs only on their jaws, crocodiles have similar organs on almost every scale on their bodies. The function of the DPRs on the jaws is clear; to catch prey, but it is still not clear what the function is of the organs on the rest of the body. The receptors flatten when exposed to increased osmotic pressure, such as that experienced when swimming in sea water hyperosmotic to the body fluids. When contact between the integument and the surrounding sea water solution is blocked, crocodiles are found to lose their ability to discriminate salinities. It has been proposed that the flattening of the sensory organ in hyperosmotic sea water is sensed by the animal as "touch", but interpreted as chemical information about its surroundings. This might be why in alligators they are absent on the rest of the body. Hunting and diet Crocodiles are ambush predators, waiting for fish or land animals to come close, then rushing out to attack. Crocodiles mostly eat fish, amphibians, crustaceans, molluscs, birds, reptiles, and mammals, and they occasionally cannibalize smaller crocodiles. What a crocodile eats varies greatly with species, size and age. From the mostly fish-eating species, like the slender-snouted and freshwater crocodiles, to the larger species like the Nile crocodile and the saltwater crocodile that prey on large mammals, such as buffalo, deer and wild boar, diet shows great diversity. Diet is also greatly affected by the size and age of the individual within the same species. All young crocodiles hunt mostly invertebrates and small fish, gradually moving on to larger prey. Being ectothermic (cold-blooded) predators, they have a very slow metabolism, so they can survive long periods without food. Despite their appearance of being slow, crocodiles have a very fast strike and are top predators in their environment, and various species have been observed attacking and killing other predators such as sharks and big cats. Crocodiles are also known to be aggressive scavengers who feed upon carrion and steal from other predators. Evidence suggests that crocodiles also feed upon fruits, based on the discovery of seeds in stools and stomachs from many subjects as well as accounts of them feeding. Crocodiles have the most acidic stomach of any vertebrate. They can easily digest bones, hooves and horns. The BBC TV reported that a Nile crocodile that has lurked a long time underwater to catch prey builds up a large oxygen debt. When it has caught and eaten that prey, it closes its right aortic arch and uses its left aortic arch to flush blood loaded with carbon dioxide from its muscles directly to its stomach; the resulting excess acidity in its blood supply makes it much easier for the stomach lining to secrete more stomach acid to quickly dissolve bulks of swallowed prey flesh and bone. Many large crocodilians swallow stones (called gastroliths or stomach stones), which may act as ballast to balance their bodies or assist in crushing food, similar to grit ingested by birds. Herodotus claimed that Nile crocodiles had a symbiotic relationship with certain birds, such as the Egyptian plover, which enter the crocodile's mouth and pick leeches feeding on the crocodile's blood; with no evidence of this interaction actually occurring in any crocodile species, it is most likely mythical or allegorical fiction. Bite Since they feed by grabbing and holding onto their prey, they have evolved sharp teeth for piercing and holding onto flesh, and powerful muscles to close the jaws and hold them shut. The teeth are not well-suited to tearing flesh off of large prey items as are the dentition and claws of many mammalian carnivores, the hooked bills and talons of raptorial birds, or the serrated teeth of sharks. However, this is an advantage rather than a disadvantage to the crocodile since the properties of the teeth allow it to hold onto prey with the least possibility of the prey animal escaping. Cutting teeth, combined with the exceptionally high bite force, would pass through flesh easily enough to leave an escape opportunity for prey. The jaws can bite down with immense force, by far the strongest bite of any animal. The force of a large crocodile's bite is more than , which was measured in a Nile crocodile, in the field; comparing to for a Rottweiler, for a hyena, for an American alligator, and for the largest confirmed great white shark. A long saltwater crocodile has been confirmed as having the strongest bite force ever recorded for an animal in a laboratory setting. It was able to apply a bite force value of , and thus surpassed the previous record of made by a long American alligator. Taking the measurements of several crocodiles as reference, the bite forces of 6-m individuals were estimated at . The study, led by Dr. Gregory M. Erickson, also shed light on the larger, extinct species of crocodilians. Since crocodile anatomy has changed only slightly over the last 80 million years, current data on modern crocodilians can be used to estimate the bite force of extinct species. An Deinosuchus would apply a force of , nearly twice that of the latest, higher bite force estimations of Tyrannosaurus (). The extraordinary bite of crocodilians is a result of their anatomy. The space for the jaw muscle in the skull is very large, which is easily visible from the outside as a bulge at each side. The muscle is so stiff, it is almost as hard as bone to touch, as if it were the continuum of the skull. Another trait is that most of the muscle in a crocodile's jaw is arranged for clamping down. Despite the strong muscles to close the jaw, crocodiles have extremely small and weak muscles to open the jaw. Crocodiles can thus be subdued for study or transport by taping their jaws or holding their jaws shut with large rubber bands cut from automobile inner tubes. Locomotion Crocodiles can move quickly over short distances, even out of water. The land speed record for a crocodile is measured in a galloping Australian freshwater crocodile. Maximum speed varies between species. Some species can gallop, including Cuban crocodiles, Johnston's crocodiles, New Guinea crocodiles, African dwarf crocodiles, and even small Nile crocodiles. The fastest means by which most species can move is a "belly run", in which the body moves in a snake-like (sinusoidal) fashion, limbs splayed out to either side paddling away frantically while the tail whips to and fro. Crocodiles can reach speeds of when they "belly run", and often faster if slipping down muddy riverbanks. When a crocodile walks quickly, it holds its legs in a straighter and more upright position under its body, which is called the "high walk". This walk allows a speed of up to 5 km/h. Crocodiles may possess a homing instinct. In northern Australia, three rogue saltwater crocodiles were relocated by helicopter, but returned to their original locations within three weeks, based on data obtained from tracking devices attached to them. Longevity Measuring crocodile age is unreliable, although several techniques are used to derive a reasonable guess. The most common method is to measure lamellar growth rings in bones and teeth—each ring corresponds to a change in growth rate which typically occurs once a year between dry and wet seasons. Bearing these inaccuracies in mind, it can be safely said that all crocodile species have an average lifespan of at least 30–40 years, and in the case of larger species an average of 60–70 years. The oldest crocodiles appear to be the largest species. C. porosus is estimated to live around 70 years on average, with limited evidence of some individuals exceeding 100 years. In captivity, some individuals are claimed to have lived for over a century. A male crocodile lived to an estimated age of 110–115 years in a Russian zoo in Yekaterinburg. Named Kolya, he joined the zoo around 1913 to 1915, fully grown, after touring in an animal show, and lived until 1995. A male freshwater crocodile lived to an estimated age of 120–140 years at the Australia Zoo. Known affectionately as "Mr. Freshie", he was rescued around 1970 by Bob Irwin and Steve Irwin, after being shot twice by hunters and losing an eye as a result, and lived until 2010. Crocworld Conservation Centre, in Scottburgh, South Africa, claims to have a male Nile crocodile that was born in 1900. Named Henry, the crocodile is said to have lived in Botswana along the Okavango River, according to centre director Martin Rodrigues. Social behaviour and vocalization Crocodiles are the most social of reptiles. Even though they do not form social groups, many species congregate in certain sections of rivers, tolerating each other at times of feeding and basking. Most species are not highly territorial, with the exception of the saltwater crocodile, which is a highly territorial and aggressive species: a mature, male saltwater crocodile will not tolerate any other males at any time of the year, but most other species are more flexible. There is a certain form of hierarchy in crocodiles: the largest and heaviest males are at the top, having access to the best basking site, while females are priority during a group feeding of a big kill or carcass. A good example of the hierarchy in crocodiles would be the case of the Nile crocodile. This species clearly displays all of these behaviours. Studies in this area are not thorough, however, and many species are yet to be studied in greater detail. Mugger crocodiles are also known to show toleration in group feedings and tend to congregate in certain areas. However, males of all species are aggressive towards each other during mating season, to gain access to females. Crocodiles are also the most vocal of all reptiles, producing a wide variety of sounds during various situations and conditions, depending on species, age, size and sex. Depending on the context, some species can communicate over 20 different messages through vocalizations alone. Some of these vocalizations are made during social communication, especially during territorial displays towards the same sex and courtship with the opposite sex; the common concern being reproduction. Therefore most conspecific vocalization is made during the breeding season, with the exception being year-round territorial behaviour in some species and quarrels during feeding. Crocodiles also produce different distress calls and in aggressive displays to their own kind and other animals; notably other predators during interspecific predatory confrontations over carcasses and terrestrial kills. Specific vocalisations include — Chirp: When about to hatch, the young make a "peeping" noise, which encourages the female to excavate the nest. The female then gathers the hatchlings in her mouth and transports them to the water, where they remain in a group for several months, protected by the female Distress call: A high-pitched call used mostly by younger animals to alert other crocodiles to imminent danger or an animal being attacked. Threat call: A hissing sound that has also been described as a coughing noise. Hatching call: Emitted by a female when breeding to alert other crocodiles that she has laid eggs in her nest. Bellowing: Male crocodiles are especially vociferous. Bellowing choruses occur most often in the spring when breeding groups congregate, but can occur at any time of year. To bellow, males noticeably inflate as they raise the tail and head out of water, slowly waving the tail back and forth. They then puff out the throat and with a closed mouth, begin to vibrate air. Just before bellowing, males project an infrasonic signal at about 10 Hz through the water, which vibrates the ground and nearby objects. These low-frequency vibrations travel great distances through both air and water to advertise the male's presence and are so powerful they result in the water's appearing to "dance". Reproduction Crocodiles lay eggs, which are laid in either holes or mound nests, depending on species. A hole nest is usually excavated in sand and a mound nest is usually constructed out of vegetation. Nesting periods range from a few weeks up to six months. Courtship takes place in a series of behavioural interactions that include a variety of snout rubbing and submissive display that can take a long time. Mating always takes place in water, where the pair can be observed mating several times. Females can build or dig several trial nests which appear incomplete and abandoned later. Egg-laying usually takes place at night and about 30–40 minutes. Females are highly protective of their nests and young. The eggs are hard shelled, but translucent at the time of egg-laying. Depending on the species of crocodile, 7 to 95 eggs are laid. Crocodile embryos do not have sex chromosomes, and unlike humans, sex is not determined genetically. Sex is determined by temperature, where at or less most hatchlings are females and at , offspring are of both sexes. A temperature of gives mostly males whereas above in some species continues to give males, but in other species resulting in females, which are sometimes called high-temperature females. Temperature also affects growth and survival rate of the young, which may explain the sexual dimorphism in crocodiles. The average incubation period is around 80 days, and also is dependent on temperature and species that usually ranges from 65 to 95 days. The eggshell structure is very conservative through evolution but there are enough changes to tell different species apart by their eggshell microstructure. Scutes may play a role in calcium storage for eggshell formation. At the time of hatching, the young start calling within the eggs. They have an egg-tooth at the tip of their snouts, which is developed from the skin, and that helps them pierce out of the shell. Hearing the calls, the female usually excavates the nest and sometimes takes the unhatched eggs in her mouth, slowly rolling the eggs to help the process. The young is usually carried to the water in the mouth. She would then introduce her hatchlings to the water and even feed them. The mother would then take care of her young for over a year before the next mating season. In the absence of the mother crocodile, the father would act in her place to take care of the young. However, even with a sophisticated parental nurturing, young crocodiles have a very high mortality rate due to their vulnerability to predation. A group of hatchlings is called a pod or crèche and may be protected for months. Cognition Crocodiles possess some advanced cognitive abilities. They can observe and use patterns of prey behaviour, such as when prey come to the river to drink at the same time each day. Vladimir Dinets of the University of Tennessee, observed that crocodiles use twigs as bait for birds looking for nesting material. They place sticks on their snouts and partly submerge themselves. When the birds swooped in to get the sticks, the crocodiles then catch the birds. Crocodiles only do this in spring nesting seasons of the birds, when there is high demand for sticks to be used for building nests. Vladimir also discovered other similar observations from various scientists, some dating back to the 19th century. Aside from using sticks, crocodiles are also capable of cooperative hunting. Large numbers of crocodiles swim in circles to trap fish and take turns snatching them. In hunting larger prey, crocodiles swarm in, with one holding the prey down as the others rip it apart. According to a 2015 study, crocodiles engage in all three main types of play behaviour recorded in animals: locomotor play, play with objects and social play. Play with objects is reported most often, but locomotor play such as repeatedly sliding down slopes, and social play such as riding on the backs of other crocodiles is also reported. This behaviour was exhibited with conspecifics and mammals and is apparently not uncommon, though has been difficult to observe and interpret in the past due to obvious dangers of interacting with large carnivores. Relationship with humans Danger to humans The larger species of crocodiles are very dangerous to humans, mainly because of their ability to strike before the person can react. The saltwater crocodile and Nile crocodile are the most dangerous, killing hundreds of people each year in parts of Southeast Asia and Africa. The mugger crocodile and American crocodile are also dangerous to humans. Crocodile products Crocodiles are protected in many parts of the world, but are also farmed commercially. Their hides are tanned and used to make leather goods such as shoes and handbags; crocodile meat is also considered a delicacy. The most commonly farmed species are the saltwater and Nile crocodiles, while a hybrid of the saltwater and the rare Siamese crocodile is also bred in Asian farms. Farming has resulted in an increase in the saltwater crocodile population in Australia, as eggs are usually harvested from the wild, so landowners have an incentive to conserve their habitat. Crocodile leather can be made into goods such as wallets, briefcases, purses, handbags, belts, hats, and shoes. Crocodile oil has been used for various purposes. Crocodiles were eaten by Vietnamese while they were taboo and off limits for Chinese. Vietnamese women who married Chinese men adopted the Chinese taboo. Crocodile meat is consumed in some countries, such as Australia, Ethiopia, Thailand, South Africa, China, and Cuba (in pickled form). It is also occasionally eaten as an "exotic" delicacy in the western world. Cuts of meat include backstrap and tail fillet. Due to high demand for crocodile products, TRAFFIC states that 1,418,487 Nile Crocodile skins were exported from Africa between 2006 and 2015. Crocodile hunting and conservation Aboriginal Australians harvested eggs and hunted crocodiles in a sustainable way for many thousands of years. The Brinkin people (aka Marrithiyal) of the Daly River in the Northern Territory (NT) used harpoons and bamboo, and even their own hands to capture crocodiles for food. After settlement of northern Australia, in the late-19th and early 20th centuries, non-Indigenous people killed individual crocodiles, mostly by locals to protect the population, or novelty-seeking visitors, or just opportunistically, so numbers were not noticeably reduced. From the 1930s, commercial hunting began, with Aboriginal people often employed to kill the crocodiles using traditional methods. From the 1940s to the 1960s, hunting began on a larger scale using .303 rifles. They were hunted for leather, with the skins shipped to plants in capital cities. Western Australia banned hunting freshwater crocodiles in 1962 and saltwater crocodiles in 1970, while NT bans were brought in 1964 and 1971; Queensland did not pass such legislation. The federal government later banned the export of crocodile skins, which brought commercial hunting to an end in Queensland. They have been a protected species since the 1970s, when numbers were down to approximately 3,000 in the NT at the lowest estimate. In 2021, after several attacks on humans by the "salties" and an estimated population of around 200,000 had been reached, Queensland politician Bob Katter called for the reintroduction of hunting. In religion and mythology Crocodiles have appeared in various forms in religions across the world. Ancient Egypt had Sobek, the crocodile-headed god, with his cult-city Crocodilopolis, as well as Taweret, the goddess of childbirth and fertility, with the back and tail of a crocodile. The Jukun shrine in the Wukari Federation, Nigeria is dedicated to crocodiles in thanks for their aid during migration. In Madagascar various peoples such as the Sakalava and Antandroy see crocodiles as ancestor spirits and under local fady often offer them food; in the case of the latter at least a crocodile features prominently as an ancestor deity. Crocodiles appear in different forms in Hinduism. Varuna, a Vedic and Hindu god, rides a part-crocodile makara; his consort Varuni rides a crocodile. Similarly the goddess personifications of the Ganga and Yamuna rivers are often depicted as riding crocodiles. Also in India, in Goa, crocodile worship is practised, including the annual Mannge Thapnee ceremony. Sikh warriors known as nihang also have connections with crocodiles. Nihang may come from the Persian word for a mythical sea creature (). The term owes its origin to Mughal historians, who compared the ferocity of the Akali with that of crocodiles. The meaning of Akali in Sikhism however, is the immortal army of Akal (god). In Latin America, Cipactli was the giant earth crocodile of the Aztec and other Nahua peoples. The name of Surabaya, Indonesia, is locally believed to be derived from the words "suro" (shark) and "boyo" (crocodile), two creatures which, in a local myth, fought each other in order to gain the title of "the strongest and most powerful animal" in the area. It was said that the two powerful animals agreed for a truce and set boundaries; that the shark's domain would be in the sea while the crocodile's domain would be on the land. However one day the shark swam into the river estuary to hunt, this angered the crocodile, who declared it his territory. The Shark argued that the river was a water-realm which meant that it was shark territory, while the crocodile argued that the river flowed deep inland, so it was therefore crocodile territory. A ferocious fight resumed as the two animals bit each other. Finally the shark was badly bitten and fled to the open sea, and the crocodile finally ruled the estuarine area that today is the city. Another source alludes to a Jayabaya prophecy—a 12th-century psychic king of Kediri Kingdom—as he foresaw a fight between a giant white shark and a giant white crocodile taking place in the area, which is sometimes interpreted as a foretelling of the Mongol invasion of Java, a major conflict between the forces of the Kublai Khan, Mongol ruler of China, and those of Raden Wijaya's Majapahit in 1293. The two animals are now used as the city's symbol, with the two facing and circling each other, as depicted in a statue appropriately located near the entrance to the city zoo (see photo on the Surabaya page). In language and as symbols The term "crocodile tears" (and equivalents in other languages) refers to a false, insincere display of emotion, such as a hypocrite crying fake tears of grief. It is derived from an ancient anecdote that crocodiles weep in order to lure their prey, or that they cry for the victims they are eating, first told in the Bibliotheca by Photios I of Constantinople. The story is repeated in bestiaries such as De bestiis et aliis rebus. This tale was first spread widely in English in the stories of the Travels of Sir John Mandeville in the 14th century, and appears in several of Shakespeare's plays. In fact, crocodiles can and do generate tears, but they do not actually cry. In the UK, a row of schoolchildren walking in pairs, or two by two is known as "crocodile". Fashion logos The French clothing company Lacoste features a crocodile in its logo. The American shoe company Crocs also uses this imagery in its logo.	Biology and health sciences	Reptiles	null
37895	https://en.wikipedia.org/wiki/Collatz%20conjecture	Collatz conjecture	The Collatz conjecture is one of the most famous unsolved problems in mathematics. The conjecture asks whether repeating two simple arithmetic operations will eventually transform every positive integer into 1. It concerns sequences of integers in which each term is obtained from the previous term as follows: if a term is even, the next term is one half of it. If a term is odd, the next term is 3 times the previous term plus 1. The conjecture is that these sequences always reach 1, no matter which positive integer is chosen to start the sequence. The conjecture has been shown to hold for all positive integers up to , but no general proof has been found. It is named after the mathematician Lothar Collatz, who introduced the idea in 1937, two years after receiving his doctorate. The sequence of numbers involved is sometimes referred to as the hailstone sequence, hailstone numbers or hailstone numerals (because the values are usually subject to multiple descents and ascents like hailstones in a cloud), or as wondrous numbers. Paul Erdős said about the Collatz conjecture: "Mathematics may not be ready for such problems." Jeffrey Lagarias stated in 2010 that the Collatz conjecture "is an extraordinarily difficult problem, completely out of reach of present day mathematics". However, though the Collatz conjecture itself remains open, efforts to solve the problem have led to new techniques and many partial results. Statement of the problem Consider the following operation on an arbitrary positive integer: If the number is even, divide it by two. If the number is odd, triple it and add one. In modular arithmetic notation, define the function as follows: Now form a sequence by performing this operation repeatedly, beginning with any positive integer, and taking the result at each step as the input at the next. In notation: (that is: is the value of applied to recursively times; ). The Collatz conjecture is: This process will eventually reach the number 1, regardless of which positive integer is chosen initially. That is, for each , there is some with . If the conjecture is false, it can only be because there is some starting number which gives rise to a sequence that does not contain 1. Such a sequence would either enter a repeating cycle that excludes 1, or increase without bound. No such sequence has been found. The smallest such that is called the stopping time of . Similarly, the smallest such that is called the total stopping time of . If one of the indexes or doesn't exist, we say that the stopping time or the total stopping time, respectively, is infinite. The Collatz conjecture asserts that the total stopping time of every is finite. It is also equivalent to saying that every has a finite stopping time. Since is even whenever is odd, one may instead use the "shortcut" form of the Collatz function: This definition yields smaller values for the stopping time and total stopping time without changing the overall dynamics of the process. Empirical data For instance, starting with and applying the function without "shortcut", one gets the sequence . The number takes longer to reach 1: . The sequence for , listed and graphed below, takes 111 steps (41 steps through odd numbers, in bold), climbing as high as 9232 before descending to 1. Numbers with a total stopping time longer than that of any smaller starting value form a sequence beginning with: 1, 2, 3, 6, 7, 9, 18, 25, 27, 54, 73, 97, 129, 171, 231, 313, 327, 649, 703, 871, 1161, 2223, 2463, 2919, 3711, 6171, ... . The starting values whose maximum trajectory point is greater than that of any smaller starting value are as follows: 1, 2, 3, 7, 15, 27, 255, 447, 639, 703, 1819, 4255, 4591, 9663, 20895, 26623, 31911, 60975, 77671, 113383, 138367, 159487, 270271, 665215, 704511, ... Number of steps for to reach 1 are 0, 1, 7, 2, 5, 8, 16, 3, 19, 6, 14, 9, 9, 17, 17, 4, 12, 20, 20, 7, 7, 15, 15, 10, 23, 10, 111, 18, 18, 18, 106, 5, 26, 13, 13, 21, 21, 21, 34, 8, 109, 8, 29, 16, 16, 16, 104, 11, 24, 24, ... The starting value having the largest total stopping time while being less than 10 is 9, which has 19 steps, less than 100 is 97, which has 118 steps, less than 1000 is 871, which has 178 steps, less than 104 is 6171, which has 261 steps, less than 105 is , which has 350 steps, less than 106 is , which has 524 steps, less than 107 is , which has 685 steps, less than 108 is , which has 949 steps, less than 109 is , which has 986 steps, less than 1010 is , which has 1132 steps, less than 1011 is , which has 1228 steps, less than 1012 is , which has 1348 steps. These numbers are the lowest ones with the indicated step count, but not necessarily the only ones below the given limit. As an example, has 1132 steps, as does . The starting values having the smallest total stopping time with respect to their number of digits (in base 2) are the powers of two since is halved times to reach 1, and is never increased. Visualizations Supporting arguments Although the conjecture has not been proven, most mathematicians who have looked into the problem think the conjecture is true because experimental evidence and heuristic arguments support it. Experimental evidence The conjecture has been checked by computer for all starting values up to 268 ≈ . All values tested so far converge to 1. This computer evidence is still not rigorous proof that the conjecture is true for all starting values, as counterexamples may be found when considering very large (or possibly immense) positive integers, as in the case of the disproven Pólya conjecture and Mertens conjecture. However, such verifications may have other implications. Certain constraints on any non-trivial cycle, such as lower bounds on the length of the cycle, can be proven based on the value of the lowest term in the cycle. Therefore, computer searches to rule out cycles that have a small lowest term can strengthen these constraints. A probabilistic heuristic If one considers only the odd numbers in the sequence generated by the Collatz process, then each odd number is on average of the previous one. (More precisely, the geometric mean of the ratios of outcomes is .) This yields a heuristic argument that every Hailstone sequence should decrease in the long run, although this is not evidence against other cycles, only against divergence. The argument is not a proof because it assumes that Hailstone sequences are assembled from uncorrelated probabilistic events. (It does rigorously establish that the 2-adic extension of the Collatz process has two division steps for every multiplication step for almost all 2-adic starting values.) Stopping times As proven by Riho Terras, almost every positive integer has a finite stopping time. In other words, almost every Collatz sequence reaches a point that is strictly below its initial value. The proof is based on the distribution of parity vectors and uses the central limit theorem. In 2019, Terence Tao improved this result by showing, using logarithmic density, that almost all (in the sense of logarithmic density) Collatz orbits are descending below any given function of the starting point, provided that this function diverges to infinity, no matter how slowly. Responding to this work, Quanta Magazine wrote that Tao "came away with one of the most significant results on the Collatz conjecture in decades". Lower bounds In a computer-aided proof, Krasikov and Lagarias showed that the number of integers in the interval that eventually reach 1 is at least equal to for all sufficiently large . Cycles In this part, consider the shortcut form of the Collatz function A cycle is a sequence of distinct positive integers where , , ..., and . The only known cycle is of period 2, called the trivial cycle. Cycle length The length of a non-trivial cycle is known to be at least (or without shortcut). If it can be shown that for all positive integers less than the Collatz sequences reach 1, then this bound would raise to ( without shortcut). In fact, Eliahou (1993) proved that the period of any non-trivial cycle is of the form where , and are non-negative integers, and . This result is based on the simple continued fraction expansion of . -cycles A -cycle is a cycle that can be partitioned into contiguous subsequences, each consisting of an increasing sequence of odd numbers, followed by a decreasing sequence of even numbers. For instance, if the cycle consists of a single increasing sequence of odd numbers followed by a decreasing sequence of even numbers, it is called a 1-cycle. Steiner (1977) proved that there is no 1-cycle other than the trivial . Simons (2005) used Steiner's method to prove that there is no 2-cycle. Simons and de Weger (2005) extended this proof up to 68-cycles; there is no -cycle up to . Hercher extended the method further and proved that there exists no k-cycle with . As exhaustive computer searches continue, larger values may be ruled out. To state the argument more intuitively; we do not have to search for cycles that have less than 92 subsequences, where each subsequence consists of consecutive ups followed by consecutive downs. Other formulations of the conjecture In reverse There is another approach to prove the conjecture, which considers the bottom-up method of growing the so-called Collatz graph. The Collatz graph is a graph defined by the inverse relation So, instead of proving that all positive integers eventually lead to 1, we can try to prove that 1 leads backwards to all positive integers. For any integer , if and only if . Equivalently, if and only if . Conjecturally, this inverse relation forms a tree except for the 1–2–4 loop (the inverse of the 4–2–1 loop of the unaltered function defined in the Statement of the problem section of this article). When the relation of the function is replaced by the common substitute "shortcut" relation , the Collatz graph is defined by the inverse relation, For any integer , if and only if . Equivalently, if and only if . Conjecturally, this inverse relation forms a tree except for a 1–2 loop (the inverse of the 1–2 loop of the function f(n) revised as indicated above). Alternatively, replace the with where and is the highest power of 2 that divides (with no remainder). The resulting function maps from odd numbers to odd numbers. Now suppose that for some odd number , applying this operation times yields the number 1 (that is, ). Then in binary, the number can be written as the concatenation of strings where each is a finite and contiguous extract from the representation of . The representation of therefore holds the repetends of , where each repetend is optionally rotated and then replicated up to a finite number of bits. It is only in binary that this occurs. Conjecturally, every binary string that ends with a '1' can be reached by a representation of this form (where we may add or delete leading '0's to ). As an abstract machine that computes in base two Repeated applications of the Collatz function can be represented as an abstract machine that handles strings of bits. The machine will perform the following three steps on any odd number until only one remains: Append to the (right) end of the number in binary (giving ); Add this to the original number by binary addition (giving ); Remove all trailing s (that is, repeatedly divide by 2 until the result is odd). Example The starting number 7 is written in base two as . The resulting Collatz sequence is: 111 1111 10110 10111 100010 100011 110100 11011 101000 1011 10000 As a parity sequence For this section, consider the shortcut form of the Collatz function If is the parity of a number, that is and , then we can define the Collatz parity sequence (or parity vector) for a number as , where , and . Which operation is performed, or , depends on the parity. The parity sequence is the same as the sequence of operations. Using this form for , it can be shown that the parity sequences for two numbers and will agree in the first terms if and only if and are equivalent modulo . This implies that every number is uniquely identified by its parity sequence, and moreover that if there are multiple Hailstone cycles, then their corresponding parity cycles must be different. Applying the function times to the number will give the result , where is the result of applying the function times to , and is how many increases were encountered during that sequence. For example, for there are 3 increases as 1 iterates to 2, 1, 2, 1, and finally to 2 so the result is ; for there is only 1 increase as 1 rises to 2 and falls to 1 so the result is . When is then there will be rises and the result will be . The power of 3 multiplying is independent of the value of ; it depends only on the behavior of . This allows one to predict that certain forms of numbers will always lead to a smaller number after a certain number of iterations: for example, becomes after two applications of and becomes after four applications of . Whether those smaller numbers continue to 1, however, depends on the value of . As a tag system For the Collatz function in the shortcut form Hailstone sequences can be computed by the 2-tag system with production rules , , . In this system, the positive integer is represented by a string of copies of , and iteration of the tag operation halts on any word of length less than 2. (Adapted from De Mol.) The Collatz conjecture equivalently states that this tag system, with an arbitrary finite string of as the initial word, eventually halts (see Tag system for a worked example). Extensions to larger domains Iterating on all integers An extension to the Collatz conjecture is to include all integers, not just positive integers. Leaving aside the cycle 0 → 0 which cannot be entered from outside, there are a total of four known cycles, which all nonzero integers seem to eventually fall into under iteration of . These cycles are listed here, starting with the well-known cycle for positive : Odd values are listed in large bold. Each cycle is listed with its member of least absolute value (which is always odd) first. The generalized Collatz conjecture is the assertion that every integer, under iteration by , eventually falls into one of the four cycles above or the cycle 0 → 0. Iterating on rationals with odd denominators The Collatz map can be extended to (positive or negative) rational numbers which have odd denominators when written in lowest terms. The number is taken to be 'odd' or 'even' according to whether its numerator is odd or even. Then the formula for the map is exactly the same as when the domain is the integers: an 'even' such rational is divided by 2; an 'odd' such rational is multiplied by 3 and then 1 is added. A closely related fact is that the Collatz map extends to the ring of 2-adic integers, which contains the ring of rationals with odd denominators as a subring. When using the "shortcut" definition of the Collatz map, it is known that any periodic parity sequence is generated by exactly one rational. Conversely, it is conjectured that every rational with an odd denominator has an eventually cyclic parity sequence (Periodicity Conjecture). If a parity cycle has length and includes odd numbers exactly times at indices , then the unique rational which generates immediately and periodically this parity cycle is For example, the parity cycle has length 7 and four odd terms at indices 0, 2, 3, and 6. It is repeatedly generated by the fraction as the latter leads to the rational cycle Any cyclic permutation of is associated to one of the above fractions. For instance, the cycle is produced by the fraction For a one-to-one correspondence, a parity cycle should be irreducible, that is, not partitionable into identical sub-cycles. As an illustration of this, the parity cycle and its sub-cycle are associated to the same fraction when reduced to lowest terms. In this context, assuming the validity of the Collatz conjecture implies that and are the only parity cycles generated by positive whole numbers (1 and 2, respectively). If the odd denominator of a rational is not a multiple of 3, then all the iterates have the same denominator and the sequence of numerators can be obtained by applying the " " generalization of the Collatz function 2-adic extension The function is well-defined on the ring of 2-adic integers, where it is continuous and measure-preserving with respect to the 2-adic measure. Moreover, its dynamics is known to be ergodic. Define the parity vector function acting on as The function is a 2-adic isometry. Consequently, every infinite parity sequence occurs for exactly one 2-adic integer, so that almost all trajectories are acyclic in . An equivalent formulation of the Collatz conjecture is that Iterating on real or complex numbers The Collatz map can be extended to the real line by choosing any function which evaluates to when is an even integer, and to either or (for the "shortcut" version) when is an odd integer. This is called an interpolating function. A simple way to do this is to pick two functions and , where: and use them as switches for our desired values: . One such choice is and . The iterations of this map lead to a dynamical system, further investigated by Marc Chamberland. He showed that the conjecture does not hold for positive real numbers since there are infinitely many fixed points, as well as orbits escaping monotonically to infinity. The function has two attracting cycles of period : and . Moreover, the set of unbounded orbits is conjectured to be of measure . Letherman, Schleicher, and Wood extended the study to the complex plane. They used Chamberland's function for complex sine and cosine and added the extra term , where is any entire function. Since this expression evaluates to zero for real integers, the extended function is an interpolation of the Collatz map to the complex plane. The reason for adding the extra term is to make all integers critical points of . With this, they show that no integer is in a Baker domain, which implies that any integer is either eventually periodic or belongs to a wandering domain. They conjectured that the latter is not the case, which would make all integer orbits finite. Most of the points have orbits that diverge to infinity. Coloring these points based on how fast they diverge produces the image on the left, for . The inner black regions and the outer region are the Fatou components, and the boundary between them is the Julia set of , which forms a fractal pattern, sometimes called a "Collatz fractal". There are many other ways to define a complex interpolating function, such as using the complex exponential instead of sine and cosine: , which exhibit different dynamics. In this case, for instance, if , then . The corresponding Julia set, shown on the right, consists of uncountably many curves, called hairs, or rays. Optimizations Time–space tradeoff The section As a parity sequence above gives a way to speed up simulation of the sequence. To jump ahead steps on each iteration (using the function from that section), break up the current number into two parts, (the least significant bits, interpreted as an integer), and (the rest of the bits as an integer). The result of jumping ahead is given by . The values of (or better ) and can be precalculated for all possible -bit numbers , where is the result of applying the function times to , and is the number of odd numbers encountered on the way. For example, if , one can jump ahead 5 steps on each iteration by separating out the 5 least significant bits of a number and using (0...31, 5) = { 0, 3, 2, 2, 2, 2, 2, 4, 1, 4, 1, 3, 2, 2, 3, 4, 1, 2, 3, 3, 1, 1, 3, 3, 2, 3, 2, 4, 3, 3, 4, 5 }, (0...31, 5) = { 0, 2, 1, 1, 2, 2, 2, 20, 1, 26, 1, 10, 4, 4, 13, 40, 2, 5, 17, 17, 2, 2, 20, 20, 8, 22, 8, 71, 26, 26, 80, 242 }. This requires precomputation and storage to speed up the resulting calculation by a factor of , a space–time tradeoff. Modular restrictions For the special purpose of searching for a counterexample to the Collatz conjecture, this precomputation leads to an even more important acceleration, used by Tomás Oliveira e Silva in his computational confirmations of the Collatz conjecture up to large values of . If, for some given and , the inequality holds for all , then the first counterexample, if it exists, cannot be modulo . For instance, the first counterexample must be odd because , smaller than ; and it must be 3 mod 4 because , smaller than . For each starting value which is not a counterexample to the Collatz conjecture, there is a for which such an inequality holds, so checking the Collatz conjecture for one starting value is as good as checking an entire congruence class. As increases, the search only needs to check those residues that are not eliminated by lower values of . Only an exponentially small fraction of the residues survive. For example, the only surviving residues mod 32 are 7, 15, 27, and 31. Integers divisible by 3 cannot form a cycle, so these integers do not need to be checked as counter examples. Syracuse function If is an odd integer, then is even, so with odd and . The Syracuse function is the function from the set of positive odd integers into itself, for which . Some properties of the Syracuse function are: For all , . (Because .) In more generality: For all and odd , . (Here is function iteration notation.) For all odd , The Collatz conjecture is equivalent to the statement that, for all in , there exists an integer such that . Undecidable generalizations In 1972, John Horton Conway proved that a natural generalization of the Collatz problem is algorithmically undecidable. Specifically, he considered functions of the form where are rational numbers which are so chosen that is always an integer. The standard Collatz function is given by , , , , . Conway proved that the problem Given and , does the sequence of iterates reach ? is undecidable, by representing the halting problem in this way. Closer to the Collatz problem is the following universally quantified problem: Given , does the sequence of iterates reach , for all ? Modifying the condition in this way can make a problem either harder or easier to solve (intuitively, it is harder to justify a positive answer but might be easier to justify a negative one). Kurtz and Simon proved that the universally quantified problem is, in fact, undecidable and even higher in the arithmetical hierarchy; specifically, it is -complete. This hardness result holds even if one restricts the class of functions by fixing the modulus to 6480. Iterations of in a simplified version of this form, with all equal to zero, are formalized in an esoteric programming language called FRACTRAN. In computational complexity Collatz and related conjectures are often used when studying computation complexity. The connection is made through the Busy Beaver function, where BB(n) is the maximum number of steps taken by any n state Turing machine that halts. There is a 15 state Turing machine that halts if and only if a conjecture by Paul Erdős (closely related to the Collatz conjecture) is false. Hence if BB(15) was known, and this machine did not stop in that number of steps, it would be known to run forever and hence no counterexamples exist (which proves the conjecture true). This is a completely impractical way to settle the conjecture; instead it is used to suggest that BB(15) will be very hard to compute, at least as difficult as settling this Collatz-like conjecture.	Mathematics	Sums and products	null
37910	https://en.wikipedia.org/wiki/Spacecraft	Spacecraft	A spacecraft is a vehicle that is designed to fly and operate in outer space. Spacecraft are used for a variety of purposes, including communications, Earth observation, meteorology, navigation, space colonization, planetary exploration, and transportation of humans and cargo. All spacecraft except single-stage-to-orbit vehicles cannot get into space on their own, and require a launch vehicle (carrier rocket). On a sub-orbital spaceflight, a space vehicle enters space and then returns to the surface without having gained sufficient energy or velocity to make a full Earth orbit. For orbital spaceflights, spacecraft enter closed orbits around the Earth or around other celestial bodies. Spacecraft used for human spaceflight carry people on board as crew or passengers from start or on orbit (space stations) only, whereas those used for robotic space missions operate either autonomously or telerobotically. Robotic spacecraft used to support scientific research are space probes. Robotic spacecraft that remain in orbit around a planetary body are artificial satellites. To date, only a handful of interstellar probes, such as Pioneer 10 and 11, Voyager 1 and 2, and New Horizons,are on trajectories that leave the Solar System. Orbital spacecraft may be recoverable or not. Most are not. Recoverable spacecraft may be subdivided by a method of reentry to Earth into non-winged space capsules and winged spaceplanes. Recoverable spacecraft may be reusable (can be launched again or several times, like the SpaceX Dragon and the Space Shuttle orbiters) or expendable (like the Soyuz). In recent years, more space agencies are tending towards reusable spacecraft. Humanity has achieved space flight, but only a few nations have the technology for orbital launches: Russia (Roscosmos), the United States (NASA), the member states of the European Space Agency, Japan (JAXA), China (CNSA), India (ISRO), Taiwan (TSA), Israel (ISA), Iran (ISA), and North Korea (NADA). In addition, several private companies have developed or are developing the technology for orbital launches independently from government agencies. Two prominent examples of such companies are SpaceX and Blue Origin. History A German V-2 became the first spacecraft when it reached an altitude of 189 km in June 1944 in Peenemünde, Germany. Sputnik 1 was the first artificial satellite. It was launched into an elliptical low Earth orbit (LEO) by the Soviet Union on 4 October 1957. The launch ushered in new political, military, technological, and scientific developments; while the Sputnik launch was a single event, it marked the start of the Space Age. Apart from its value as a technological first, Sputnik 1 also helped to identify the upper atmospheric layer's density, by measuring the satellite's orbital changes. It also provided data on radio-signal distribution in the ionosphere. Pressurized nitrogen in the satellite's false body provided the first opportunity for meteoroid detection. Sputnik 1 was launched during the International Geophysical Year from Site No.1/5, at the 5th Tyuratam range, in Kazakh SSR (now at the Baikonur Cosmodrome). The satellite travelled at , taking 96.2 minutes to complete an orbit, and emitted radio signals at 20.005 and 40.002 MHz While Sputnik 1 was the first spacecraft to orbit the Earth, other human-made objects had previously reached an altitude of 100 km, which is the height required by the international organization Fédération Aéronautique Internationale to count as a spaceflight. This altitude is called the Kármán line. In particular, in the 1940s there were several test launches of the V-2 rocket, some of which reached altitudes well over 100 km. Crewed and uncrewed spacecraft Crewed spacecraft As of 2016, only three nations have flown crewed spacecraft: USSR/Russia, USA, and China. The first crewed spacecraft was Vostok 1, which carried Soviet cosmonaut Yuri Gagarin into space in 1961, and completed a full Earth orbit. There were five other crewed missions which used a Vostok spacecraft. The second crewed spacecraft was named Freedom 7, and it performed a sub-orbital spaceflight in 1961 carrying American astronaut Alan Shepard to an altitude of just over . There were five other crewed missions using Mercury spacecraft. Other Soviet crewed spacecraft include the Voskhod, Soyuz, flown uncrewed as Zond/L1, L3, TKS, and the Salyut and Mir crewed space stations. Other American crewed spacecraft include the Gemini spacecraft, the Apollo spacecraft including the Apollo Lunar Module, the Skylab space station, the Space Shuttle with undetached European Spacelab and private US Spacehab space stations-modules, and the SpaceX Crew Dragon configuration of their Dragon 2. US company Boeing also developed and flown a spacecraft of their own, the CST-100, commonly referred to as Starliner, but a crewed flight is yet to occur. China developed, but did not fly Shuguang, and is currently using Shenzhou (its first crewed mission was in 2003). Except for the Space Shuttle and the Buran spaceplane of the Soviet Union, the latter of which only ever had one uncrewed test flight, all of the recoverable crewed orbital spacecraft were space capsules. The International Space Station, crewed since November 2000, is a joint venture between Russia, the United States, Canada and several other countries. Uncrewed spacecraft Uncrewed spacecraft are spacecraft without people on board. Uncrewed spacecraft may have varying levels of autonomy from human input; they may be remote controlled, remote guided or even autonomous, meaning they have a pre-programmed list of operations, which they will execute unless otherwise instructed. Many space missions are more suited to telerobotic rather than crewed operation, due to lower cost and lower risk factors. In addition, some planetary destinations such as Venus or the vicinity of Jupiter are too hostile for human survival. Outer planets such as Saturn, Uranus, and Neptune are too distant to reach with current crewed spaceflight technology, so telerobotic probes are the only way to explore them. Telerobotics also allows exploration of regions that are vulnerable to contamination by Earth micro-organisms since spacecraft can be sterilized. Humans can not be sterilized in the same way as a spaceship, as they coexist with numerous micro-organisms, and these micro-organisms are also hard to contain within a spaceship or spacesuit. Multiple space probes were sent to study Moon, the planets, the Sun, multiple small Solar System bodies (comets and asteroids). Special class of uncrewed spacecraft is space telescopes, a telescope in outer space used to observe astronomical objects. The first operational telescopes were the American Orbiting Astronomical Observatory, OAO-2 launched in 1968, and the Soviet Orion 1 ultraviolet telescope aboard space station Salyut 1 in 1971. Space telescopes avoid the filtering and distortion (scintillation) of electromagnetic radiation which they observe, and avoid light pollution which ground-based observatories encounter. The best-known examples are Hubble Space Telescope and James Webb Space Telescope. Cargo spacecraft are designed to carry cargo, possibly to support space stations' operation by transporting food, propellant and other supplies. Automated cargo spacecraft have been used since 1978 and have serviced Salyut 6, Salyut 7, Mir, the International Space Station and Tiangong space station. Other Some spacecrafts can operate as both a crewed and uncrewed spacecraft. For example, the Buran spaceplane could operate autonomously but also had manual controls, though it never flew with crew onboard. Other dual crewed/uncrewed spacecrafts include: SpaceX Dragon 2, Dream Chaser, and Tianzhou. Types of spacecraft Communications satellite A communications satellite is an artificial satellite that relays and amplifies radio telecommunication signals via a transponder; it creates a communication channel between a source transmitter and a receiver at different locations on Earth. Communications satellites are used for television, telephone, radio, internet, and military applications. Many communications satellites are in geostationary orbit above the equator, so that the satellite appears stationary at the same point in the sky; therefore the satellite dish antennas of ground stations can be aimed permanently at that spot and do not have to move to track the satellite. Others form satellite constellations in low Earth orbit, where antennas on the ground have to follow the position of the satellites and switch between satellites frequently. The high frequency radio waves used for telecommunications links travel by line of sight and so are obstructed by the curve of the Earth. The purpose of communications satellites is to relay the signal around the curve of the Earth allowing communication between widely separated geographical points. Communications satellites use a wide range of radio and microwave frequencies. To avoid signal interference, international organizations have regulations for which frequency ranges or "bands" certain organizations are allowed to use. This allocation of bands minimizes the risk of signal interference. Cargo spacecraft Cargo or resupply spacecraft are robotic spacecraft that are designed specifically to carry cargo, possibly to support space stations' operation by transporting food, propellant and other supplies. Automated cargo spacecraft have been used since 1978 and have serviced Salyut 6, Salyut 7, Mir, the International Space Station and Tiangong space station. As of 2023, three different cargo spacecraft are used to supply the International Space Station: Russian Progress, American SpaceX Dragon 2 and Cygnus. Chinese Tianzhou is used to supply Tiangong space station. Space probes Space probes are robotic spacecraft that are sent to explore deep space, or astronomical bodies other than Earth. They are distinguished from landers by the fact that they work in open space, not on planetary surfaces or in planetary atmospheres. Being robotic eliminates the need for expensive, heavy life support systems (the Apollo crewed Moon landings required the use of the Saturn V rocket that cost over a billion dollars per launch, adjusted for inflation) and so allows for lighter, less expensive rockets. Space probes have visited every planet in the Solar System and Pluto, and the Parker Solar Probe has an orbit that, at its closest point, is in the Sun's chromosphere. There are five space probes that are escaping the Solar System, these are Voyager 1, Voyager 2, Pioneer 10, Pioneer 11, and New Horizons. Voyager program The identical Voyager probes, weighing , were launched in 1977 to take advantage of a rare alignment of Jupiter, Saturn, Uranus and Neptune that would allow a spacecraft to visit all four planets in one mission, and get to each destination faster by using gravity assist. In fact, the rocket that launched the probes (the Titan IIIE) could not even send the probes to the orbit of Saturn, yet Voyager 1 is travelling at roughly and Voyager 2 moves at about kilometres per second as of 2023. In 2012, Voyager 1 exited the heliosphere, followed by Voyager 2 in 2018. Voyager 1 actually launched 16 days after Voyager 2 but it reached Jupiter sooner because Voyager 2 was taking a longer route that allowed it to visit Uranus and Neptune, whereas Voyager 1 did not visit Uranus or Neptune, instead choosing to fly past Saturn’s satellite Titan. As of August 2023, Voyager 1 has passed 160 astronomical units, which means it is over 160 times farther from the Sun than Earth is. This makes it the farthest spacecraft from the Sun. Voyager 2 is 134 AU away from the Sun as of August 2023. NASA provides real time data of their distances and data from the probe’s cosmic ray detectors. Because of the probe’s declining power output and degradation of the RTGs over time, NASA has had to shut down certain instruments to conserve power. The probes may still have some scientific instruments on until the mid-2020s or perhaps the 2030s. After 2036, they will both be out of range of the Deep Space Network. Space telescopes A space telescope or space observatory is a telescope in outer space used to observe astronomical objects. Space telescopes avoid the filtering and distortion of electromagnetic radiation which they observe, and avoid light pollution which ground-based observatories encounter. They are divided into two types: satellites which map the entire sky (astronomical survey), and satellites which focus on selected astronomical objects or parts of the sky and beyond. Space telescopes are distinct from Earth imaging satellites, which point toward Earth for satellite imaging, applied for weather analysis, espionage, and other types of information gathering. Landers A lander is a type of spacecraft that makes a soft landing on the surface of an astronomical body other than Earth. Some landers, such as Philae and the Apollo Lunar Module, land entirely by using their fuel supply, however many landers (and landings of spacecraft on Earth) use aerobraking, especially for more distant destinations. This involves the spacecraft using a fuel burn to change its trajectory so it will pass through a planet (or a moon's) atmosphere. Drag caused by the spacecraft hitting the atmosphere enables it to slow down without using fuel, however this generates very high temperatures and so adds a requirement for a heat shield of some sort. Space capsules Space capsules are a type of spacecraft that can return from space at least once. They have a blunt shape, do not usually contain much more fuel than needed, and they do not possess wings unlike spaceplanes. They are the simplest form of recoverable spacecraft, and so the most commonly used. The first such capsule was the Vostok capsule built by the Soviet Union, that carried the first person in space, Yuri Gagarin. Other examples include the Soyuz and Orion capsules, built by the Soviet Union and NASA, respectively. Spaceplanes Spaceplanes are spacecraft that are built in the shape of, and function as, airplanes. The first example of such was the North American X-15 spaceplane, which conducted two crewed flights which reached an altitude of over in the 1960s. This first reusable spacecraft was air-launched on a suborbital trajectory on July 19, 1963. The first reusable orbital spaceplane was the Space Shuttle orbiter. The first orbiter to fly in space, the Space Shuttle Columbia, was launched by the USA on the 20th anniversary of Yuri Gagarin's flight, on April 12, 1981. During the Shuttle era, six orbiters were built, all of which have flown in the atmosphere and five of which have flown in space. Enterprise was used only for approach and landing tests, launching from the back of a Boeing 747 SCA and gliding to deadstick landings at Edwards AFB, California. The first Space Shuttle to fly into space was Columbia, followed by Challenger, Discovery, Atlantis, and Endeavour. Endeavour was built to replace Challenger when it was lost in January 1986. Columbia broke up during reentry in February 2003. The first autonomous reusable spaceplane was the Buran-class shuttle, launched by the USSR on November 15, 1988, although it made only one flight and this was uncrewed. This spaceplane was designed for a crew and strongly resembled the U.S. Space Shuttle, although its drop-off boosters used liquid propellants and its main engines were located at the base of what would be the external tank in the American Shuttle. Lack of funding, complicated by the dissolution of the USSR, prevented any further flights of Buran. The Space Shuttle was subsequently modified to allow for autonomous re-entry in case of necessity. Per the Vision for Space Exploration, the Space Shuttle was retired in 2011 mainly due to its old age and high cost of program reaching over a billion dollars per flight. The Shuttle's human transport role is to be replaced by SpaceX's SpaceX Dragon 2 and Boeing's CST-100 Starliner. Dragon 2's first crewed flight occurred on May 30, 2020. The Shuttle's heavy cargo transport role is to be replaced by expendable rockets such as the Space Launch System and ULA's Vulcan rocket, as well as the commercial launch vehicles. Scaled Composites' SpaceShipOne was a reusable suborbital spaceplane that carried pilots Mike Melvill and Brian Binnie on consecutive flights in 2004 to win the Ansari X Prize. The Spaceship Company built a successor SpaceShipTwo. A fleet of SpaceShipTwos operated by Virgin Galactic was planned to begin reusable private spaceflight carrying paying passengers in 2014, but was delayed after the crash of VSS Enterprise. Space Shuttle The Space Shuttle is a retired reusable Low Earth Orbital launch system. It consisted of two reusable Solid Rocket Boosters that landed by parachute, were recovered at sea, and were the most powerful rocket motors ever made until they were superseded by those of NASA’s SLS rocket, with a liftoff thrust of , which soon increased to per booster, and were fueled by a combination of PBAN and APCP, the Space Shuttle Orbiter, with 3 RS-25 engines that used a liquid oxygen/liquid hydrogen propellant combination, and the bright orange throwaway Space Shuttle external tank from which the RS-25 engines sourced their fuel. The orbiter was a spaceplane that was launched at NASA’s Kennedy Space Centre and landed mainly at the Shuttle Landing Facility, which is part of Kennedy Space Centre. A second launch site, Vandenberg Space Launch Complex 6 in California, was revamped so it could be used to launch the shuttles, but it was never used. The launch system could lift about into an eastward Low Earth Orbit. Each orbiter weighed roughly , however the different orbiters had differing weights and thus payloads, with Columbia being the heaviest orbiter, Challenger being lighter than Columbia but still heavier than the other three. The orbiter structure was mostly composed of aluminium alloy. The orbiter had seven seats for crew members, though on STS-61-A the launch took place with 8 crew onboard. The orbiters had wide by long payload bays and also were equipped with a CanadaArm1, an upgraded version of which is used on the International Space Station. The heat shield (or Thermal Protection System) of the orbiter, used to protect it from extreme levels of heat during atmospheric reentry and the cold of space, was made up of different materials depending on weight and how much heating a particular area on the shuttle would receive during reentry, which ranged from over to under . The orbiter was manually operated, though an autonomous landing system was added while the shuttle was still on service. It had an in orbit maneouvreing system known as the Orbital Manoeuvring System, which used the hypergolic propellants monomethylhydrazine (MMH) and dinitrogen tetroxide, which was used for orbital insertion, changes to orbits and the deorbit burn. Though the shuttle’s goals were to drastically decrease launch costs, it did not do so, ending up being much more expensive than similar expendable launchers. This was due to expensive refurbishment costs and the external tank being expended. Once a landing had occurred, the SRBs and many parts of the orbiter had to be disassembled for inspection, which was long and arduous. Furthermore, the RS-25 engines had to be replaced every few flights. Each of the heat shielding tiles had to go in one specific area on the orbiter, increasing complexity more. Adding to this, the shuttle was a rather dangerous system, with fragile heat shielding tiles, some being so fragile that one could easily scrape it off by hand, often having been damaged in many flights. After 30 years in service from 1981 to 2011 and 135 flights, the shuttle was retired from service due to the cost of maintaining the shuttles, and the 3 remaining orbiters (the other two were destroyed in accidents) were prepared to be displayed in museums. Other Some spacecraft do not fit particularly well into any of the general spacecraft categories. This is a list of these spacecraft. SpaceX Starship Starship is a spacecraft and second stage under development by American aerospace company SpaceX. Stacked atop its booster, Super Heavy, it composes the identically named Starship super heavy-lift space vehicle. The spacecraft is designed to transport both crew and cargo to a variety of destinations, including Earth orbit, the Moon, Mars, and potentially beyond. It is intended to enable long duration interplanetary flights for a crew of up to 100 people. It will also be capable of point-to-point transport on Earth, enabling travel to anywhere in the world in less than an hour. Furthermore, the spacecraft will be used to refuel other Starship vehicles to allow them to reach higher orbits to and other space destinations. Elon Musk, the CEO of SpaceX, estimated in a tweet that 8 launches would be needed to completely refuel a Starship in low Earth orbit, extrapolating this from Starship's payload to orbit and how much fuel a fully fueled Starship contains. To land on bodies without an atmosphere, such as the Moon, Starship will fire its engines and thrusters to slow down. Mission Extension Vehicle The Mission Extension Vehicle is a robotic spacecraft designed to prolong the life on another spacecraft. It works by docking to its target spacecraft, then correcting its orientation or orbit. This also allows it to rescue a satellite which is in the wrong orbit by using its own fuel to move its target to the correct orbit. The project is currently managed by Northrop Grumman Innovation Systems. As of 2023, 2 have been launched. The first launched on a Proton rocket on 9 October 2019, and did a rendezvous with Intelsat-901 on 25 February 2020. It will remain with the satellite until 2025 before the satellite is moved to a final graveyard orbit and the vehicle does a rendezvous with another satellite. The other one launched on an Ariane 5 rocket on 15 August 2020. Subsystems A spacecraft astrionics system comprises different subsystems, depending on the mission profile. Spacecraft subsystems are mounted in the satellite bus and may include attitude determination and control (variously called ADAC, ADC, or ACS), guidance, navigation and control (GNC or GN&C), communications (comms), command and data handling (CDH or C&DH), power (EPS), thermal control (TCS), propulsion, and structures. Attached to the bus are typically payloads. Life support Spacecraft intended for human spaceflight must also include a life support system for the crew. Attitude control A spacecraft needs an attitude control subsystem to be correctly oriented in space and respond to external torques and forces properly. This may use reaction wheels or it may use small rocket thrusters. The altitude control subsystem consists of sensors and actuators, together with controlling algorithms. The attitude-control subsystem permits proper pointing for the science objective, sun pointing for power to the solar arrays and earth pointing for communications. GNC Guidance refers to the calculation of the commands (usually done by the CDH subsystem) needed to steer the spacecraft where it is desired to be. Navigation means determining a spacecraft's orbital elements or position. Control means adjusting the path of the spacecraft to meet mission requirements. Command and data handling The C&DH subsystem receives commands from the communications subsystem, performs validation and decoding of the commands, and distributes the commands to the appropriate spacecraft subsystems and components. The CDH also receives housekeeping data and science data from the other spacecraft subsystems and components, and packages the data for storage on a data recorder or transmission to the ground via the communications subsystem. Other functions of the CDH include maintaining the spacecraft clock and state-of-health monitoring. Communications Spacecraft, both robotic and crewed, have various communications systems for communication with terrestrial stations and for inter-satellite service. Technologies include space radio station and optical communication. In addition, some spacecraft payloads are explicitly for the purpose of ground–ground communication using receiver/retransmitter electronic technologies. Power Spacecraft need an electrical power generation and distribution subsystem for powering the various spacecraft subsystems. For spacecraft near the Sun, solar panels are frequently used to generate electrical power. Spacecraft designed to operate in more distant locations, for example Jupiter, might employ a radioisotope thermoelectric generator (RTG) to generate electrical power. Electrical power is sent through power conditioning equipment before it passes through a power distribution unit over an electrical bus to other spacecraft components. Batteries are typically connected to the bus via a battery charge regulator, and the batteries are used to provide electrical power during periods when primary power is not available, for example when a low Earth orbit spacecraft is eclipsed by Earth. Thermal control Spacecraft must be engineered to withstand transit through Earth's atmosphere and the space environment. They must operate in a vacuum with temperatures potentially ranging across hundreds of degrees Celsius as well as (if subject to reentry) in the presence of plasmas. Material requirements are such that either high melting temperature, low density materials such as beryllium and reinforced carbon–carbon or (possibly due to the lower thickness requirements despite its high density) tungsten or ablative carbon–carbon composites are used. Depending on mission profile, spacecraft may also need to operate on the surface of another planetary body. The thermal control subsystem can be passive, dependent on the selection of materials with specific radiative properties. Active thermal control makes use of electrical heaters and certain actuators such as louvers to control temperature ranges of equipments within specific ranges. Spacecraft propulsion Spacecraft may or may not have a propulsion subsystem, depending on whether or not the mission profile calls for propulsion. The Swift spacecraft is an example of a spacecraft that does not have a propulsion subsystem. Typically though, LEO spacecraft include a propulsion subsystem for altitude adjustments (drag make-up maneuvers) and inclination adjustment maneuvers. A propulsion system is also needed for spacecraft that perform momentum management maneuvers. Components of a conventional propulsion subsystem include fuel, tankage, valves, pipes, and thrusters. The thermal control system interfaces with the propulsion subsystem by monitoring the temperature of those components, and by preheating tanks and thrusters in preparation for a spacecraft maneuver. Structures Spacecraft must be engineered to withstand launch loads imparted by the launch vehicle, and must have a point of attachment for all the other subsystems. Depending on mission profile, the structural subsystem might need to withstand loads imparted by entry into the atmosphere of another planetary body, and landing on the surface of another planetary body. Payload The payload depends on the mission of the spacecraft, and is typically regarded as the part of the spacecraft "that pays the bills". Typical payloads could include scientific instruments (cameras, telescopes, or particle detectors, for example), cargo, or a human crew. Ground segment The ground segment, though not technically part of the spacecraft, is vital to the operation of the spacecraft. Typical components of a ground segment in use during normal operations include a mission operations facility where the flight operations team conducts the operations of the spacecraft, a data processing and storage facility, ground stations to radiate signals to and receive signals from the spacecraft, and a voice and data communications network to connect all mission elements. Launch vehicle The launch vehicle propels the spacecraft from Earth's surface, through the atmosphere, and into an orbit, the exact orbit being dependent on the mission configuration. The launch vehicle may be expendable or reusable. In a single stage to orbit rocket, the rocket can be considered a spacecraft itself. Spacecraft records Fastest spacecraft Parker Solar Probe (estimated at first sun close pass, will reach at final perihelion) Helios I and II Solar Probes () Furthest spacecraft from the Sun Voyager 1 at 156.13 AU as of April 2022, traveling outward at about Pioneer 10 at 122.48 AU as of December 2018, traveling outward at about Voyager 2 at 122.82 AU as of January 2020, traveling outward at about Pioneer 11 at 101.17 AU as of December 2018, traveling outward at about	Technology	Basics_10	null
37913	https://en.wikipedia.org/wiki/Escape%20velocity	Escape velocity	In celestial mechanics, escape velocity or escape speed is the minimum speed needed for an object to escape from contact with or orbit of a primary body, assuming: Ballistic trajectory – no other forces are acting on the object, such as propulsion and friction No other gravity-producing objects exist. Although the term escape velocity is common, it is more accurately described as a speed than as a velocity because it is independent of direction. Because gravitational force between two objects depends on their combined mass, the escape speed also depends on mass. For artificial satellites and small natural objects, the mass of the object makes a negligible contribution to the combined mass, and so is often ignored. Escape speed varies with distance from the center of the primary body, as does the velocity of an object traveling under the gravitational influence of the primary. If an object is in a circular or elliptical orbit, its speed is always less than the escape speed at its current distance. In contrast if it is on a hyperbolic trajectory its speed will always be higher than the escape speed at its current distance. (It will slow down as it gets to greater distance, but do so asymptotically approaching a positive speed.) An object on a parabolic trajectory will always be traveling exactly the escape speed at its current distance. It has precisely balanced positive kinetic energy and negative gravitational potential energy; it will always be slowing down, asymptotically approaching zero speed, but never quite stop. Escape velocity calculations are typically used to determine whether an object will remain in the gravitational sphere of influence of a given body. For example, in solar system exploration it is useful to know whether a probe will continue to orbit the Earth or escape to a heliocentric orbit. It is also useful to know how much a probe will need to slow down in order to be gravitationally captured by its destination body. Rockets do not have to reach escape velocity in a single maneuver, and objects can also use a gravity assist to siphon kinetic energy away from large bodies. Precise trajectory calculations require taking into account small forces like atmospheric drag, radiation pressure, and solar wind. A rocket under continuous or intermittent thrust (or an object climbing a space elevator) can attain escape at any non-zero speed, but the minimum amount of energy required to do so is always the same. Calculation Escape speed at a distance d from the center of a spherically symmetric primary body (such as a star or a planet) with mass M is given by the formula where: G is the universal gravitational constant () g = GM/d2 is the local gravitational acceleration (or the surface gravity, when ). The value GM is called the standard gravitational parameter, or μ, and is often known more accurately than either G or M separately. When given an initial speed V greater than the escape speed v, the object will asymptotically approach the hyperbolic excess speed v, satisfying the equation: For example, with the definitional value for standard gravity of , the escape velocity is . Energy required For an object of mass the energy required to escape the Earth's gravitational field is GMm / r, a function of the object's mass (where r is radius of the Earth, nominally 6,371 kilometres (3,959 mi), G is the gravitational constant, and M is the mass of the Earth, ). A related quantity is the specific orbital energy which is essentially the sum of the kinetic and potential energy divided by the mass. An object has reached escape velocity when the specific orbital energy is greater than or equal to zero. Conservation of energy The existence of escape velocity can be thought of as a consequence of conservation of energy and an energy field of finite depth. For an object with a given total energy, which is moving subject to conservative forces (such as a static gravity field) it is only possible for the object to reach combinations of locations and speeds which have that total energy; places which have a higher potential energy than this cannot be reached at all. Adding speed (kinetic energy) to an object expands the region of locations it can reach, until, with enough energy, everywhere to infinity becomes accessible. The formula for escape velocity can be derived from the principle of conservation of energy. For the sake of simplicity, unless stated otherwise, we assume that an object will escape the gravitational field of a uniform spherical planet by moving away from it and that the only significant force acting on the moving object is the planet's gravity. Imagine that a spaceship of mass m is initially at a distance r from the center of mass of the planet, whose mass is M, and its initial speed is equal to its escape velocity, v. At its final state, it will be an infinite distance away from the planet, and its speed will be negligibly small. Kinetic energy K and gravitational potential energy Ug are the only types of energy that we will deal with (we will ignore the drag of the atmosphere), so by the conservation of energy, We can set Kfinal = 0 because final velocity is arbitrarily small, and = 0 because final gravitational potential energy is defined to be zero a long distance away from a planet, so Relativistic The same result is obtained by a relativistic calculation, in which case the variable r represents the radial coordinate or reduced circumference of the Schwarzschild metric. Scenarios From the surface of a body An alternative expression for the escape velocity v particularly useful at the surface on the body is: where r is the distance between the center of the body and the point at which escape velocity is being calculated and g is the gravitational acceleration at that distance (i.e., the surface gravity). For a body with a spherically symmetric distribution of mass, the escape velocity v from the surface is proportional to the radius assuming constant density, and proportional to the square root of the average density ρ. where . This escape velocity is relative to a non-rotating frame of reference, not relative to the moving surface of the planet or moon, as explained below. From a rotating body The escape velocity relative to the surface of a rotating body depends on direction in which the escaping body travels. For example, as the Earth's rotational velocity is 465 m/s at the equator, a rocket launched tangentially from the Earth's equator to the east requires an initial velocity of about 10.735 km/s relative to the moving surface at the point of launch to escape whereas a rocket launched tangentially from the Earth's equator to the west requires an initial velocity of about 11.665 km/s relative to that moving surface. The surface velocity decreases with the cosine of the geographic latitude, so space launch facilities are often located as close to the equator as feasible, e.g. the American Cape Canaveral (latitude 28°28′ N) and the French Guiana Space Centre (latitude 5°14′ N). Practical considerations In most situations it is impractical to achieve escape velocity almost instantly, because of the acceleration implied, and also because if there is an atmosphere, the hypersonic speeds involved (on Earth a speed of 11.2 km/s, or 40,320 km/h) would cause most objects to burn up due to aerodynamic heating or be torn apart by atmospheric drag. For an actual escape orbit, a spacecraft will accelerate steadily out of the atmosphere until it reaches the escape velocity appropriate for its altitude (which will be less than on the surface). In many cases, the spacecraft may be first placed in a parking orbit (e.g. a low Earth orbit at 160–2,000 km) and then accelerated to the escape velocity at that altitude, which will be slightly lower (about 11.0 km/s at a low Earth orbit of 200 km). The required additional change in speed, however, is far less because the spacecraft already has a significant orbital speed (in low Earth orbit speed is approximately 7.8 km/s, or 28,080 km/h). From an orbiting body The escape velocity at a given height is times the speed in a circular orbit at the same height, (compare this with the velocity equation in circular orbit). This corresponds to the fact that the potential energy with respect to infinity of an object in such an orbit is minus two times its kinetic energy, while to escape the sum of potential and kinetic energy needs to be at least zero. The velocity corresponding to the circular orbit is sometimes called the first cosmic velocity, whereas in this context the escape velocity is referred to as the second cosmic velocity. For a body in an elliptical orbit wishing to accelerate to an escape orbit the required speed will vary, and will be greatest at periapsis when the body is closest to the central body. However, the orbital speed of the body will also be at its highest at this point, and the change in velocity required will be at its lowest, as explained by the Oberth effect. Barycentric escape velocity Escape velocity can either be measured as relative to the other, central body or relative to center of mass or barycenter of the system of bodies. Thus for systems of two bodies, the term escape velocity can be ambiguous, but it is usually intended to mean the barycentric escape velocity of the less massive body. Escape velocity usually refers to the escape velocity of zero mass test particles. For zero mass test particles we have that the 'relative to the other' and the 'barycentric' escape velocities are the same, namely . But when we can't neglect the smaller mass (say m) we arrive at slightly different formulas. Because the system has to obey the law of conservation of momentum we see that both the larger and the smaller mass must be accelerated in the gravitational field. Relative to the center of mass the velocity of the larger mass (v, for planet) can be expressed in terms of the velocity of the smaller mass (v, for rocket). We get . The 'barycentric' escape velocity now becomes , while the 'relative to the other' escape velocity becomes . Height of lower-velocity trajectories Ignoring all factors other than the gravitational force between the body and the object, an object projected vertically at speed v from the surface of a spherical body with escape velocity v and radius R will attain a maximum height h satisfying the equation which, solving for h results in where is the ratio of the original speed v to the escape velocity v. Unlike escape velocity, the direction (vertically up) is important to achieve maximum height. Trajectory If an object attains exactly escape velocity, but is not directed straight away from the planet, then it will follow a curved path or trajectory. Although this trajectory does not form a closed shape, it can be referred to as an orbit. Assuming that gravity is the only significant force in the system, this object's speed at any point in the trajectory will be equal to the escape velocity at that point due to the conservation of energy, its total energy must always be 0, which implies that it always has escape velocity; see the derivation above. The shape of the trajectory will be a parabola whose focus is located at the center of mass of the planet. An actual escape requires a course with a trajectory that does not intersect with the planet, or its atmosphere, since this would cause the object to crash. When moving away from the source, this path is called an escape orbit. Escape orbits are known as orbits. C is the characteristic energy, −GM/2a, where a is the semi-major axis length, which is infinite for parabolic trajectories. If the body has a velocity greater than escape velocity then its path will form a hyperbolic trajectory and it will have an excess hyperbolic velocity, equivalent to the extra energy the body has. A relatively small extra delta-v above that needed to accelerate to the escape speed can result in a relatively large speed at infinity. Some orbital manoeuvres make use of this fact. For example, at a place where escape speed is 11.2 km/s, the addition of 0.4 km/s yields a hyperbolic excess speed of 3.02 km/s: If a body in circular orbit (or at the periapsis of an elliptical orbit) accelerates along its direction of travel to escape velocity, the point of acceleration will form the periapsis of the escape trajectory. The eventual direction of travel will be at 90 degrees to the direction at the point of acceleration. If the body accelerates to beyond escape velocity the eventual direction of travel will be at a smaller angle, and indicated by one of the asymptotes of the hyperbolic trajectory it is now taking. This means the timing of the acceleration is critical if the intention is to escape in a particular direction. If the speed at periapsis is , then the eccentricity of the trajectory is given by: This is valid for elliptical, parabolic, and hyperbolic trajectories. If the trajectory is hyperbolic or parabolic, it will asymptotically approach an angle from the direction at periapsis, with The speed will asymptotically approach List of escape velocities In this table, the left-hand half gives the escape velocity from the visible surface (which may be gaseous as with Jupiter for example), relative to the centre of the planet or moon (that is, not relative to its moving surface). In the right-hand half, Ve refers to the speed relative to the central body (for example the sun), whereas Vte is the speed (at the visible surface of the smaller body) relative to the smaller body (planet or moon). The last two columns will depend precisely where in orbit escape velocity is reached, as the orbits are not exactly circular (particularly Mercury and Pluto). Deriving escape velocity using calculus Let G be the gravitational constant and let M be the mass of the earth (or other gravitating body) and m be the mass of the escaping body or projectile. At a distance r from the centre of gravitation the body feels an attractive force The work needed to move the body over a small distance dr against this force is therefore given by The total work needed to move the body from the surface r0 of the gravitating body to infinity is then In order to do this work to reach infinity, the body's minimal kinetic energy at departure must match this work, so the escape velocity v0 satisfies which results in	Physical sciences	Orbital mechanics	null
37941	https://en.wikipedia.org/wiki/Sikorsky%20UH-60%20Black%20Hawk	Sikorsky UH-60 Black Hawk	The Sikorsky UH-60 Black Hawk is a four-blade, twin-engine, medium-lift utility military helicopter manufactured by Sikorsky Aircraft. Sikorsky submitted the S-70 design for the United States Army's Utility Tactical Transport Aircraft System (UTTAS) competition in 1972. The Army designated the prototype as the YUH-60A and selected the Black Hawk as the winner of the program in 1976, after a fly-off competition with the Boeing Vertol YUH-61. Named after the Native American war leader Black Hawk, the UH-60A entered service with the U.S. Army in 1979, to replace the Bell UH-1 Iroquois as the Army's tactical transport helicopter. This was followed by the fielding of electronic warfare and special operations variants of the Black Hawk. Improved UH-60L and UH-60M utility variants have also been developed. Modified versions have also been developed for the U.S. Navy, Air Force, and Coast Guard. In addition to U.S. Army use, the UH-60 family has been exported to several nations. Black Hawks have served in combat during conflicts in Grenada, Panama, Iraq, Somalia, Ukraine, the Balkans, Afghanistan, and other areas in the Middle East. Major variants include the Sikorsky SH-60 Seahawk used for naval purposes, Sikorsky HH-60 Pave Hawk for combat search and rescue, with other upgrades for various export, VIP, and special operation variants. The latest utility variant is the UH-60M. Development Initial requirement In the late 1960s, the United States Army began forming requirements for a helicopter to replace the UH-1 Iroquois, and designated the program as the Utility Tactical Transport Aircraft System (UTTAS). The Army also initiated the development of a new, common turbine engine for its helicopters that would become the General Electric T700. Based on experience in Vietnam, the Army required significant performance, survivability and reliability improvements from both UTTAS and the new powerplant. The Army released its UTTAS request for proposals (RFP) in January 1972. The RFP also included air transport requirements. Transport within the C-130 limited the UTTAS cabin height and length. The UTTAS requirements for improved reliability, survivability and lower life-cycle costs resulted in features such as dual-engines with improved hot and high altitude performance, and a modular design (reduced maintenance footprint); run-dry gearboxes; ballistically tolerant, redundant subsystems (hydraulic, electrical and flight controls); crashworthy crew (armored) and troop seats; dual-stage oleo main landing gear; ballistically tolerant, crashworthy main structure; quieter, more robust main and tail rotor systems; and a ballistically tolerant, crashworthy fuel system. Four prototypes were constructed, with the first YUH-60A flying on 17 October 1974. Prior to the delivery of the prototypes to the US Army, a preliminary evaluation was conducted in November 1975 to ensure the aircraft could be operated safely during all testing. Three of the prototypes were delivered to the Army in March 1976, for evaluation against the rival Boeing-Vertol design, the YUH-61A, and one was kept by Sikorsky for internal research. The Army selected the UH-60 for production in December 1976. Deliveries of the UH-60A to the Army began in October 1978 and the helicopter entered service in June 1979. Upgrades and variations After entering service, the helicopter was modified for new missions and roles, including mine laying and medical evacuation. An EH-60 variant was developed to conduct electronic warfare and special operations aviation developed the MH-60 variant to support its missions. Due to weight increases from the addition of mission equipment and other changes, the Army ordered the improved UH-60L in 1987. The new model incorporated all of the modifications made to the UH-60A fleet as standard design features. The UH-60L also featured more power and lifting capability with upgraded T700-GE-701C engines and an improved gearbox, both from the SH-60B Seahawk. Its external lift capacity increased by up to . The UH-60L also incorporated the SH-60B's automatic flight control system (AFCS) for better flight control with more powerful engines. Production of the L-model began in 1989. Development of the next improved variant, the UH-60M, was approved in 2001, to extend the service life of the UH-60 design into the 2020s. The UH-60M incorporates upgraded T700-GE-701D engines, improved rotor blades, and state-of-the-art electronic instrumentation, flight controls and aircraft navigation control. After the U.S. DoD approved low-rate initial production of the new variant, manufacturing began in 2006, with the first of 22 new UH-60Ms delivered in July 2006. After an initial operational evaluation, the Army approved full-rate production and a five-year contract for 1,227 helicopters in December 2007. By March 2009, 100 UH-60M helicopters had been delivered to the Army. In November 2014, the US military ordered 102 aircraft of various H-60 types, worth $1.3 billion. Following their use in the operation to kill Osama bin Laden in May 2011, it emerged that the 160th SOAR used a secret version of the UH-60 modified with low-observable technology which enabled it to evade Pakistani radar. Analysis of the tail section, the only remaining part of the aircraft which crashed during the operation, revealed extra blades on the tail rotor and other noise reduction measures, making the craft much quieter than conventional UH-60s. The aircraft appeared to include features like special high-tech materials, harsh angles, and flat surfaces found only in stealth jets. Low observable versions of the Black Hawk have been studied as far back as the mid-1970s. In September 2012, Sikorsky was awarded a Combat Tempered Platform Demonstration (CTPD) contract to further improve the Black Hawk's durability and survivability. The company is to develop new technologies such as a zero-vibration system, adaptive flight control laws, advanced fire management, a more durable main rotor, full-spectrum crashworthiness, and damage-tolerant airframe; then they are to transition them to the helicopter. Improvements to the Black Hawk are to continue until the Future Vertical Lift program is ready to replace it. In December 2014, the 101st Airborne Division began testing new resupply equipment called the Enhanced Speed Bag System (ESBS). Soldiers in the field requiring quick resupply have depended on speed bags filled with items airdropped from a UH-60. However, all systems were ad hoc with bags not made to keep objects secure from impacts, so up to half of the airdropped items would be damaged upon hitting the ground. Started in 2011, the ESBS sought to standardize the airdrop resupply method and keep up to 90 percent of supplies intact. The system includes a hands-free reusable linear brake and expendable speed line and a multipurpose cargo bag. When the bag is deployed, the brake applies friction to the rope, slowing it down enough to keep the bag oriented down on the padded base, a honeycomb and foam kit inside to dissipate energy. The ESBS better protects helicopter-dropped supplies, and allows the Black Hawk to fly higher above the ground, up from 10 feet, while travelling , limiting exposure to ground fire. Each bag can weigh and up to six can be deployed at once, dropping at . Since supplies can be delivered more accurately and the system can be automatically released on its own, the ESBS can enable autonomous resupply from unmanned helicopters. Design The UH-60 features four-blade main and tail rotors, and is powered by two General Electric T700 turboshaft engines. The main rotor is fully articulated and has elastomeric bearings in the rotor head. The tail rotor is canted and features a rigid crossbeam. The helicopter has a long, low profile shape to meet the Army's requirement for transporting aboard a C-130 Hercules, with some disassembly. It can carry 11 troops with equipment, lift of cargo internally or of cargo (for UH-60L/M) externally by sling. The Black Hawk helicopter series can perform a wide array of missions, including the tactical transport of troops, electronic warfare, and aeromedical evacuation. A VIP version known as the VH-60N is used to transport important government officials (e.g., Congress, Executive departments) with the helicopter's call sign of "Marine One" when transporting the President of the United States. In air assault operations, it can move a squad of 11 combat troops or reposition a 105 mm M119 howitzer with 30 rounds ammunition and a four-man crew in a single lift. The Black Hawk is equipped with advanced avionics and electronics for increased survivability and capability, such as the Global Positioning System. The UH-60 can be equipped with stub wings at the top of the fuselage to carry fuel tanks or various armaments. The initial stub wing system is called External Stores Support System (ESSS). It has two pylons on each wing to carry two and two tanks in total. The four fuel tanks and associated lines and valves form the external extended range fuel system (ERFS). U.S. Army UH-60s have had their ESSS modified into the crashworthy external fuel system (CEFS) configuration, replacing the older tanks with up to four total crashworthy tanks along with self-sealing fuel lines. The ESSS can also carry of armament such as rockets, missiles and gun pods. The ESSS entered service in 1986. However, it was found that the four fuel tanks obstruct the field of fire for the door guns; thus, the external tank system (ETS), carrying two fuel tanks on the stub wings, was developed. The unit cost of the H-60 models varies due to differences in specifications, equipment and quantities. For example, the unit cost of the Army's UH-60L Black Hawk is $5.9 million while the Air Force HH-60G Pave Hawk has a unit cost of $10.2 million. Operational history The UH-60 Black Hawk is in service with 35 countries as of 2024. Australia Australia bought early model UH-60 in the 1980s, and is buying a fleet of newer versions ones in the 2020s: Australia ordered fourteen S-70A-9 Black Hawks in 1986 and an additional twenty-five Black Hawks in 1987. The first US-produced Black Hawk was delivered in 1987 to the Royal Australian Air Force (RAAF). de Havilland Australia produced thirty-eight Black Hawks under license from Sikorsky in Australia delivering the first in 1988 and the last in 1991. In 1989, the RAAF's fleet of Black Hawks was transferred to the Australian Army. The Black Hawks saw operational service in Cambodia, Papua New Guinea, Indonesia, East Timor and Pakistan. In April 2009, the then-defence chief Air Chief Marshal Angus Houston, told the government not to deploy Black Hawks to Afghanistan as at the time they "lacked armor and self-defense systems", and despite an upgrade to address this underway, it was more practical to use allies' helicopters. In 2004, the government selected the Multi-Role Helicopter (MRH-90) Taipan, a variant of the NHIndustries NH90, to replace the Black Hawk even though the Department of Defence had recommended the S‐70M Black Hawk. In January 2014, the Army began retiring the fleet of 34 Black Hawks from service (five had been lost in accidents) and had planned for this to be completed by June 2018. The Chief of Army delayed the retirement of 20 Black Hawks until 2021 to enable the Army to develop a special operations role capable MRH-90. On 10 December 2021, the S-70A-9 Black Hawks were retired from service. On the same day, amid issues with the performance of the MRH-90s the government announced that they would be replaced by UH-60M Black Hawks. In January 2023, the Army announced the acquisition of 40 UH-60Ms with deliveries commencing in 2023. Brazil Brazil received four UH-60L helicopters in 1997, for the Brazilian Army peacekeeping forces. It received six UH-60Ls configured for special forces, and search and rescue uses in 2008. It ordered ten more UH-60Ls in 2009; deliveries began in March 2011. In July 2024, the MoD authorized the purchase of 12 additional UH-60Ms, in a US$451 million plan. China In December 1983, examples of the Aerospatiale AS-332 Super Puma, Bell 214ST SuperTransport and Sikorsky S-70A-5 (N3124B) were airlifted to Lhasa for testing. These demonstrations included take-offs and landings at altitudes to and en route operations to . At the end of this testing, the People's Liberation Army purchased 24 S-70C-2s, equipped with more powerful GE T700-701A engines for improved high-altitude performance. While designated as civil variants of the S-70 for export purposes, they are operated by the People's Liberation Army Aviation units. Colombia Colombia first received UH-60s from the United States in 1987. The Colombian National Police, Colombian Aerospace Force, and Colombian Army use UH-60s to transport troops and supplies to places which are difficult to access by land for counter-insurgency (COIN) operations against drug and guerrilla organizations, for search and rescue, and for medical evacuation. Colombia also operates a militarized gunship version of the UH-60, with stub wings, locally known as Arpía (). The Colombian Army became the first worldwide operator of the S-70i with Terrain Awareness and Warning Capability (HTAWS) after taking delivery of the first two units on 13 August 2013. Israel The Israeli Air Force (IAF) received 10 surplus UH-60A Black Hawks from the United States in August 1994. Named Yanshuf () by the IAF, the UH-60A began replacing Bell 212 utility helicopters. The IAF first used the UH-60s in combat during 1996 in southern Lebanon in Operation Grapes of Wrath against Hezbollah. Mexico The Mexican Air Force ordered its first two UH-60Ls in 1991 to transport special forces units, and another four in 1994. In July and August 2009, the Federal Police used UH-60s in attacks on drug traffickers. In August 2011, the Mexican Navy received three upgraded and navalized UH-60M. On 21 April 2014, the U.S. State Department approved the sale of 18 UH-60Ms to Mexico pending approval from Congress. In September 2014, Sikorsky received a $203.6 million (~$ in ) firm-fixed-price contract modification for the 18 UH-60s designated for the Mexican Air Force. Philippines 2 S-70-A5 VIP helicopters purchased 1983 and was delivered in 1984, this Blackhawk served the 250th PAW for more than 3 decades as a Presidential VVIP transport helicopter. Only 1 remains in service with the 505th Search and Rescue Group. In March 2019, the Philippines' Department of National Defense (DND) signed a contract worth US$241.4 million (~$ in ) with Lockheed Martin's Polish subsidiary PZL Mielec for 16 Sikorsky S-70i Black Hawks to the PAF. On 10 December 2020, the PAF commissioned their first batch of six S-70i Blackhawks, with the remaining 10 to be delivered in 2021. In June 2021, the air service received a second batch of five helicopters. In November 2021, the third batch of five arrived. On 22 February 2022, DND and PZL Mielec formally signed the US$624 million contract for 32 additional S-70i Black Hawks, totalling to around 48 units ordered. This will make the Philippine Air Force the largest user of S-70i Blackhawk Helicopters globally. Poland In January 2019, Poland ordered four S-70i Black Hawks with four delivered to the Polish Special Forces in December of that same year. Another four S-70i helicopters are on order with two scheduled for delivery in 2023 and two in 2024. In July 2023, Poland launched a procurement tender for S-70i Black Hawks with a goal to order approximately 32 helicopters. Slovakia In February 2015, the U.S. State Department approved a possible Foreign Military Sale of nine UH-60Ms with associated equipment and support to Slovakia and sent it to Congress for its approval. In April 2015, Slovakia's government approved the procurement of nine UH-60Ms along with training and support. In September 2015, Slovakia ordered four UH-60Ms. The first two UH-60Ms were delivered in June 2017; the Slovak Air Force had received all nine UH-60Ms by January 2020. These are to replace its old Soviet Mil Mi-17s. In 2020, the Slovak minister of defense announced Slovakia's interest in buying two more UH-60Ms. Slovak Training Academy (European Air Services / Heli Company) from Košice, a private company, operates some older UH-60As & Bs for pilot training, aerial fire fighting and sky crane operations. Sweden Sweden requested 15 UH-60M helicopters by Foreign Military Sale in September 2010. The UH-60Ms were ordered in May 2011, and deliveries began in January 2012. In March 2013, Swedish ISAF forces began using Black Hawks in Afghanistan for MEDEVAC purposes. The UH-60Ms have been fully operational since 2017. Sweden designates it the Helicopter 16 (Hkp 16). In June 2024, Sweden ordered 12 more UH-60Ms from the US. South Korea The Republic of Korea Armed Forces is also an operator and has produced about 130 aircraft under license from Korean Air since the 1990s and domestically producing and introducing the UH-60 simulator. However, the cockpit is analog compared to the digital one in the United States, but since this business started after the 1988 Olympics, there was no such thing as a glass cockpit with an LCD monitor. Currently, the majority of South Korea's UH-60s belong to the Army, including more than 30 units operated by the Special Operations Aviation Corps. Taiwan Taiwan (Republic of China) operated S-70C-1/1A after the Republic of China Air Force received ten S-70C-1A and four S-70C-1 Bluehawk helicopters in June 1986 for Search and Rescue. Four more S-70C-6s were received in April 1998. The ROC Navy received the first of ten S-70C(M)-1s in July 1990. 11 S-70C(M)-2s were received beginning April 2000. In January 2010, the US announced approval for a Foreign Military Sale of 60 UH-60Ms to Taiwan with 30 designated for the Army, 15 for the National Airborne Service Corps (including the one that crashed off Orchid Island in 2018) and 15 for the Air Force Rescue Group (including the one that crashed 2 January 2020). Turkey Turkey has operated the UH-60 during NATO deployments to Afghanistan and the Balkans. The UH-60 has also been used in counter-terror/internal security operations. The Black Hawk competed against the AgustaWestland AW149 in the Turkish General Use Helicopter Tender, to order up to 115 helicopters and produce many of them indigenously, with Turkish Aerospace Industries responsible for final integration and assembly. On 21 April 2011, Turkey announced the selection of Sikorsky's T-70. In the course of the coup d'état attempt in Turkey on 15 July 2016, eight Turkish military personnel of various ranks landed in Greece's northeastern city of Alexandroupolis on board a Black Hawk helicopter and claimed political asylum in Greece. The helicopter was returned to Turkey shortly thereafter. Ukraine In February 2023, Ukraine's Main Directorate of Intelligence (HUR) published a video showing them operating at least two UH-60s painted in Ukrainian colors. The helicopters appeared to have minimal modifications, namely the addition of two M240 7.62 mm machine guns for defensive purposes. It was confirmed that at least one of these was purchased by a third party, Ace Aeronautics, following a Czech crowdfunding effort that raised US$6 million. On 17 March 2024, Russia claimed to have shot down a UH-60 during the March 2024 western Russia incursion, claiming it was a "troop transport" carrying 20 troops into combat. However, it was revealed to be a Mil Mi-8 instead. United States The UH-60 entered service with the U.S. Army's 101st Combat Aviation Brigade of the 101st Airborne Division in June 1979. The U.S. military first used the UH-60 in combat during the invasion of Grenada in 1983, and again in the invasion of Panama in 1989. During the Gulf War in 1991, the UH-60 participated in the largest air assault mission in U.S. Army history with over 300 helicopters involved. Two UH-60s (89-26214 and 78–23015) were shot down, both on 27 February 1991, while performing Combat Search and Rescue of other downed aircrews, an F-16C pilot and the crew of a MEDEVAC UH-1H that were shot down earlier that day. In 1993, Black Hawks featured prominently in the Battle of Mogadishu in Somalia. Black Hawks also saw action in the Balkans and Haiti in the 1990s. U.S. Army UH-60s and other helicopters conducted many air assaults and other support missions during the 2003 invasion of Iraq. The UH-60 has continued to serve in operations in Afghanistan and Iraq. Customs and Border Protection Office of Air and Marine (OAM) uses the UH-60 in its operations specifically along the southwest border. The Black Hawk has been used by OAM to interdict illegal entry into the U.S. Additionally, OAM regularly uses the UH-60 in search and rescue operations. Highly modified H-60s were employed during the U.S. Special Operations mission that resulted in the death of Osama bin Laden during Operation Neptune Spear on 1 May 2011. One such MH-60 helicopter crash-landed during the operation and was destroyed by the team before it departed in the other MH-60 and a backup MH-47 Chinook with bin Laden's remains. Two MH-47s were used for the mission to refuel the two MH-60s and as backups. News media reported that the Pakistani government granted the Chinese military access to the wreckage of the crashed 'stealth' UH-60 variant in Abbottabad; Pakistan and China denied the reports, and the U.S. government did not confirm Chinese access. The U.S. Army has signalled its intent to eventually replace the UH-60, launching the Future Long-Range Assault Aircraft (FLRAA) program in 2019, with a new helicopter planned to enter service by 2030. Bell and a joint Sikorsky-Boeing team both entered competing designs. In December 2022 it was announced that the winning design was Bell’s tilt-rotor V-280 Valor, with the US Army awarding an initial contract to develop a prototype by 2025. This award does not guarantee the eventual adoption of the V-280, which would require further contracts. As an Army program, the outcome of FLRAA will not necessarily affect UH-60 variants in service with other branches of the U.S. military. Additional users The United Arab Emirates requested 14 UH-60M helicopters and associated equipment in September 2008, through Foreign Military Sale. It had received 20 UH-60Ls by November 2010. Bahrain ordered nine UH-60Ms in 2007. In December 2011, the Royal Brunei Air Force (RBAirF / TUDB) ordered twelve S-70i helicopters, which are similar to the UH-60M; four aircraft had been received by December 2013. In June 2012, the U.S. Defense Security Cooperation Agency notified Congress that Qatar requested the purchase of twelve UH-60Ms, engines, and associated equipment. The Royal Brunei Air Force had earlier bought four UH-60, but these were later sold to Malaysia. In May 2014, Croatian Defence Minister Ante Kotromanović announced the beginning of negotiations for the purchase of 15 used Black Hawks. In October 2018, the US via Ambassador Robert Kohorst announced donation of two UH-60M helicopters with associated equipment and crew training to Croatia's Ministry of Defence, to be delivered in 2020. In October 2019, the US State Dept approved the sale of two new UH-60M Blackhawks. In February 2022, the first two helicopters were delivered to Croatia. In January 2024, the State Department approved a possible Foreign Military Sale to Croatia for 8 UH-60M helicopters and related equipment and services for an estimated cost of $500 million. The U.S. government has provided $139.4 million in financial assistance for 51 percent of the funding, as a compensation for the Croatian donation of 14 Mi-8 helicopters to Ukraine. The remaining sum is be provided by Croatia's Ministry of Defence in the three-year budget period from 2025 to 2027. The Letter of Offer and Acceptance was signed in March 2024. Delivery of all 8 Black Hawks is expected in 2028. Tunisia requested 12 armed UH-60M helicopters in July 2014 through Foreign Military Sale. In August 2014, the U.S. ambassador stated that the U.S. "will soon make available" the UH-60Ms to Tunisia. The sale of 8 helicopters was approved and helicopters were delivered 2017 and 2018. In January 2015, the Malaysian Defence Minister Hishammuddin Hussein confirmed that Royal Malaysian Air Force (RMAF) is receiving S-70A Blackhawks from the Brunei government. These helicopters, believed to be four in total, were expected to be transferred to Malaysia by September with M134D miniguns added. The four Blackhawks were delivered to Royal Brunei Air Force (RBAirF / TUDB) in 1999. In 2018, Latvia requested to buy four UH-60M Black Hawks with associated equipment for an estimated cost of $200 million (~$ in ). In August 2018, the State Department approved the possible Foreign Military Sale. The Defense Security Cooperation Agency delivered the required certification notifying Congress of the possible sale. In November 2018, Latvia ordered four UH-60Ms, and received the first two in December 2022. In 2019, Lithuania announced plans to buy six UH-60M helicopters before ordering four UH-60Ms in 2020. In July 2020, the US State Department approved the possible Foreign Military Sale of six UH-60Ms and associated equipment to Lithuania for $380 million. In November 2020, Lithuania signed a contract worth $213 million for four UH-60Ms with an option to purchase two more aircraft. Preparations are almost complete including facilities and training, with deliveries expected in late 2024. In 2019, Poland ordered four S-70i helicopters for its special forces. As of 2023 there is negotiations to purchase additional S-70i helicopters. In August 2023, the Portuguese Air Force shared a photo on twitter of the first flight of one of the six UH-60s purchased from Arista Aviation Services. The Portuguese armed forces conducted its first operation flight of its UH-60 in December 2023. In December 2023, the Hellenic Army selected the UH-60Ms for a possible order of 35 aircraft and associated equipment for an estimated cost of $1.95 billion pending the deal clears Congress. This order was approved by US and Greek governments, and a contract for 35 helicopters agreed by April 2024. In Greek service it will replace aged Bell UH-1H and Agusta-Bell AB205. Greece already operates S-70B and MH-60R helicopters. Future and potential users In February 2013, the Indonesian Army announced its interest in buying UH-60 Black Hawks to modernize its weaponry. The army wants them for combating terrorism, transnational crime, and insurgency to secure the archipelago. In August 2023, Indonesian Aerospace and Lockheed Martin signed an agreement for the procurement of 24 UH-60/S-70 Blackhawks. In 2022, the Royal Air Force and British Army expects to select a helicopter for the New Medium Helicopter program to replace several existing helicopters. Sikorsky has indicated it expects its S-70M to meet the requirement to participate in this procurement selection program. Variants The UH-60 comes in many variants and modifications. The U.S. Army variants can be fitted with stub wings to carry additional fuel tanks or weapons. Variants may have different capabilities and equipment to fulfill different roles. Utility variants YUH-60A: Initial test and evaluation version for U.S. Army. First flight on 17 October 1974. Three were built. UH-60A Black Hawk: Original U.S. Army version, carrying a crew of four and up to 11 equipped troops. Equipped with T700-GE-700 engines. Produced 1977–1989. U.S. Army is equipping UH-60As with more powerful T700-GE-701D engines and also upgrading A-models to UH-60L standards. UH-60C Black Hawk: Modified version for command and control (C2) missions. CH-60E: Proposed troop transport variant for the U.S. Marine Corps. UH-60L Black Hawk: UH-60A with upgraded T700-GE-701C engines, improved durability gearbox, and updated flight control system. Produced 1989–2007. UH-60Ls are also being equipped with the GE T700-GE-701D engine. The U.S. Army Corpus Christi Army Depot is upgrading UH-60A helicopters to the UH-60L configuration. In July 2018, Sierra Nevada Corporation proposed upgrading some converted UH-60L helicopters for the U.S. Air Force's UH-1N replacement program. UH-60M Black Hawk: Improved design wide chord rotor blades, T700-GE-701D engines (max each), improved durability gearbox, Integrated Vehicle Health Management System (IVHMS) computer, and new glass cockpit. Production began in 2006. Planned to replace older U.S. Army UH-60s. UH-60M Upgrade Black Hawk: UH-60M with fly-by-wire system and Common Avionics Architecture System (CAAS) cockpit suite. Flight testing began in August 2008. UH-60V Black Hawk: Upgraded version of the UH-60L with the electronic displays (glass cockpit) of the UH-60M. Upgrades performed by Northrop Grumman featuring a centralized processor with a partitioned, modular operational flight program enabling capabilities to be added as software-only modifications. Special purpose EH-60A Black Hawk: UH-60A with modified electrical system and stations for two electronic systems mission operators. All examples of type have been converted back to standard UH-60A configuration. YEH-60B Black Hawk: UH-60A modified for special radar and avionics installations, prototype for stand-off target acquisition system. EH-60C Black Hawk: UH-60A modified with special electronics equipment and external antenna. (All examples of type have been taken back to standard UH-60A configuration.) EUH-60L (no official name assigned): UH-60L modified with additional mission electronic equipment for Army Airborne C2. EH-60L Black Hawk: EH-60A with major mission equipment upgrade. UH-60Q Black Hawk: UH-60A modified for medical evacuation. The UH-60Q is named DUSTOFF for "dedicated unhesitating service to our fighting forces". HH-60L (no official name assigned): UH-60L extensively modified with medical mission equipment. Components include an external rescue hoist, integrated patient configuration system, environmental control system, onboard oxygen system (OBOGS), and crash-worthy ambulatory seats. HH-60M Black Hawk: UH-60M with medical mission equipment (medevac version) for U.S. Army. HH-60U: USAF UH-60M version modified with an electro-optical sensor and rescue hoist. Three in use by Air Force pilots and special mission aviators since 2011. Has 85% commonality with the HH-60W. HH-60W Jolly Green II: Modified version of the UH-60M for the U.S. Air Force as a Combat Rescue Helicopter to replace HH-60G Pave Hawks with greater fuel capacity and more internal cabin space, dubbed the "60-Whiskey". Deliveries to the USAF of the HH-50W began in 2020.The 41st Rescue Squadron received the first two HH-60W helicopters on 5 November 2020. MH-60A Black Hawk: 30 UH-60As modified with additional avionics, night vision capable cockpit, FLIR, M134 door guns, internal auxiliary fuel tanks and other Special Operations mission equipment in early 1980s for U.S. Army. Equipped with T700-GE-701 engines. Variant was used by the 160th Special Operations Aviation Regiment. The MH-60As were replaced by MH-60Ls beginning in the early 1990s and passed to Army Aviation units in the Army National Guard. MH-60L Black Hawk: Special operations modification, used by the U.S. Army's 160th Special Operations Aviation Regiment ("Night Stalkers"), based on the UH-60L with T700-701C engines. It was developed as an interim version in the late 1980s pending the fielding of the MH-60K specifically designed for the 160th SOAR(A). Equipped with many of the systems used on MH-60K, including FLIR, color weather map, auxiliary fuel system, and laser rangefinder/designator. A total of 37 MH-60Ls were built and some 10 had received an in-flight refueling probe by 2003. MH-60L DAP: The Direct Action Penetrator (DAP) is a special operations modification of the baseline MH-60L, operated by the U.S. Army's 160th Special Operations Aviation Regiment. The DAP is configured as a gunship, with no troop-carrying capacity. The DAP is equipped with ESSS or ETS stub wings, each capable of carrying configurations of the M230 Chain Gun 30 mm automatic cannon, 19-shot Hydra 70 rocket pod, AGM-114 Hellfire missiles, AIM-92 Stinger air-to-air missiles, GAU-19 gun pods, and M134 minigun pods, M134D miniguns are used as door guns. MH-60K Black Hawk: Special operations modification first ordered in 1988 for use by the U.S. Army's 160th Special Operations Aviation Regiment ("Night Stalkers"). Equipped with the in-flight refueling probe, and T700-GE-701C engines. More advanced than the MH-60L, the K-model also includes an integrated avionics system (glass cockpit), AN/APQ-174B terrain-following radar, color weather map, improved weapons capability, and various defensive systems. The MH-60K can be configured either as an assault helicopter carrying troops or as a DAP gunship. MH-60M Black Hawk: Special operations version of UH-60M for U.S. Army. Equipped with in-flight refueling probe, Rockwell Collins Common Avionics Architecture System (CAAS) glass cockpit, updated sensors and defensive systems such as the AN/APQ-187 Silent Knight terrain-following radar, and more powerful YT706-GE-700 engines. All special operations Black Hawks to be modernized to MH-60M standard by 2015. Like the K-model, the MH-60M can be configured either as an assault helicopter carrying troops or as a DAP gunship. MH-60 Black Hawk stealth helicopter: One of two (known) specially modified MH-60s used in the raid on Osama bin Laden's compound in Pakistan on 1 May 2011 was damaged in a hard landing, and was subsequently destroyed by U.S. forces. Subsequent reports state that the Black Hawk destroyed was a previously unconfirmed but rumored, modification of the design with reduced noise signature and stealth technology. The modifications are said to add several hundred pounds to the base helicopter including edge alignment panels, special coatings and anti-radar treatments for the windshields. UH-60A RASCAL: NASA-modified version for the Rotorcraft-Aircrew Systems Concepts Airborne Laboratory; a US$25M program for the study of helicopter manoeuvrability in three programs, Superaugmented Controls for Agile Maneuvering Performance (SCAMP), Automated Nap-of-the-Earth (ANOE) and Rotorcraft Agility and Pilotage Improvement Demonstration (RAPID). The UH-60A RASCAL performed a fully autonomous flight on 5 November 2012. U.S. Army personnel were on board, but the flying was done by helicopter. During a two-hour flight, the Black Hawk featured terrain sensing, trajectory generation, threat avoidance, and autonomous flight control. It was fitted with a 3D-LZ laser detection and ranging (LADAR) system. The autonomous flight was performed between 200 and 400 feet. Upon landing, the onboard technology was able to pinpoint a safe landing zone, hover, and safely bring itself down. OPBH: On 11 March 2014, Sikorsky successfully conducted the first flight demonstration of their Optionally Piloted Black Hawk (OPBH), a milestone part of the company's Manned/Unmanned Resupply Aerial Lifter (MURAL) program to provide autonomous cargo delivery for the U.S. Army. The helicopter used the company's Matrix technology (software to improve features of autonomous, optionally-piloted VTOL aircraft) to perform autonomous hover and flight operations under the control of an operator using a man-portable Ground Control Station (GCS). The MURAL program is a cooperative effort between Sikorsky, the US Army Aviation Development Directorate (ADD), and the US Army Utility Helicopters Project Office (UH PO). The purpose of creating an optionally-manned Black Hawk is to make the aircraft autonomously carry out resupply missions and expeditionary operations while increasing sorties and maintaining crew rest requirements and leaving pilots to focus more on sensitive operations. VH-60D Night Hawk: VIP-configured HH-60D, used for presidential transport by USMC. T700-GE-401C engines. Variant was later redesignated VH-60N. VH-60N White Hawk "White Top": Modified UH-60A with some features from the SH-60B/F Seahawks. Is one of the VIP-configured USMC helicopter models that perform Presidential and VIP transport as Marine One. The VH-60N entered service in 1988 and nine helicopters were delivered. VH-60M Black Hawk "Gold Top": Heavily modified UH-60M used for executive transport. Members of the Joint Chiefs, Congressional leadership, and other DoD personnel are flown on these exclusively by Alpha company 12th Aviation Battalion at Fort Belvoir, Virginia. Export versions UH-60J Black Hawk: Variant for the Japanese Air Self Defense Force and Maritime Self Defense Force produced under license by Mitsubishi Heavy Industries. Also known as the S-70-12. UH-60JA Black Hawk: Variant for the Japanese Ground Self Defense Force. It is a license produced by Mitsubishi Heavy Industries. AH-60L Arpía: Export version for Colombia developed by Elbit Systems, Sikorsky, and the Colombian Aerospace Force. It is Counter-insurgency (COIN) attack version with improved electronics, firing system, FLIR, radar, light rockets and machine guns. AH-60L Battle Hawk: Export armed version unsuccessfully tendered for Australian Army project AIR87, similar to AH-60L Arpía III. Sikorsky has also offered a Battlehawk armed version for export in the form of armament kits and upgrades. Sikorsky's Armed Black Hawk demonstrator has tested a 20 mm turreted cannon, and different guided missiles. The United Arab Emirates ordered Battlehawk kits in 2011. UH-60P Black Hawk: Version for South Korean Army, based on UH-60L with some improvements. Around 150 were produced under license by Korean Air. S-70A Black Hawk: Sikorsky's designation for Black Hawk. The designation is often used for exports. S-70A-1 Desert Hawk: Export version for the Royal Saudi Land Forces. S-70A-L1 Desert Hawk: Aeromedical evacuation version for the Royal Saudi Land Forces. S-70A-5 Black Hawk: Export version for the Philippine Air Force. S-70A-6 Black Hawk: Export version for Thailand. S-70A-9 Black Hawk: Export version for Australia, assembled under licence by Hawker de Havilland. The first eight were delivered to the Royal Australian Air Force, subsequently transferred to the Australian Army; the remainder were delivered straight to the Army after rotary-wing assets were divested by the Air Force in 1989. S-70A-11 Black Hawk: Export version for the Royal Jordanian Air Force. S-70A-12 Black Hawk: Search and rescue model for the Japanese Air Self Defense Force and Maritime Self Defense Force. Also known as the UH-60J. S-70A-14 Black Hawk: Export version for Brunei. S-70A-16 Black Hawk: Engine test bed for the Rolls-Royce/Turbomeca RTM 332. S-70A-17 Black Hawk: Export version for Turkey. S-70A-18 Black Hawk: UH-60P and HH-60P for Republic of Korea Armed Forces built under license. Sikorsky/Westland S-70-19 Black Hawk: This version is built under license in the United Kingdom by Westland. Also known as the WS-70. S-70A-20 Black Hawk: VIP transport version for Thailand. S-70A-21 Black Hawk: Export version for Egypt. S-70A-22 Black Hawk: VH-60P for South Korea built under license. Used for VIP transport by the Republic of Korea Air Force. Its fuselage is tipped with white to distinguish it from normal HH-60P. S-70A-24 Black Hawk: Export version for Mexico. S-70A-26 Black Hawk: Export version for Morocco. S-70A-27 Black Hawk: Export version for Royal Hong Kong Auxiliary Air Force and Hong Kong Government Flying Service; three built. S-70A-28D Black Hawk: Export version for Turkish Army. S-70A-30 Black Hawk: Export version for Argentine Air Force, used as a VIP transport helicopter by the Presidential fleet; one built. S-70A-33 Black Hawk: Export version for Royal Brunei Air Force. S-70A-39 Black Hawk: VIP transport version for Chile; one built. S-70A-42 Black Hawk: Export version for Austria. S-70A-43 Black Hawk: Export version for Royal Thai Army. S-70A-50 Black Hawk: Export version for Israel; 15 built. S-70C-2 Black Hawk: Export version for the People's Republic of China; 24 built. S-70i Black Hawk: International military version assembled by Sikorsky's subsidiary, PZL Mielec in Poland. S-70M Black Hawk: Modified military version assembled by Sikorsky's subsidiary, PZL Mielec in Poland from 2021. See: Sikorsky SH-60 Seahawk, Sikorsky HH-60 Pave Hawk, Piasecki X-49, and Sikorsky HH-60 Jayhawk for other Sikorsky S-70 variants. Operators See SH-60 Seahawk, HH-60 Pave Hawk, and HH-60 Jayhawk for operators of military H-60/S-70 variants; see Sikorsky S-70 for non-military operators of other H-60/S-70 family helicopters. Taliban (captured from the Afghan Air Force in August 2021) Some damaged helicopters have been repaired. Albanian Air Force - 2 (4 on order) Australian Army - 39 S-70A-9 in original orders in 1986 and 1987. Retired in 2021 with 5 lost. 40 UH-60M ordered in 2023, with 10 delivered by the end of 2024. Royal Australian Navy (see H-60 Seahawk) Austrian Air Force Royal Bahraini Air Force Brazilian Air Force Brazilian Army Brazilian Navy (see SH-60) Royal Brunei Air Force S-70i Chilean Air Force People's Liberation Army Colombian Aerospace Force AH-60L Arpía (24) Colombian Army S-70i (7 as of 2013) Croatian Air Force - 8 UH-60Ms being procured with 4 received as of March 2024. 8 more on order. Egyptian Air Force Hellenic Army - 35 UH-60Ms ordered in 2024 Indonesian Army - 22 S-70M Black Hawks on order as of 2023. Israeli Air Force Japan Air Self-Defence Force UH-60J Japan Ground Self-Defence Force UH-60JA Japan Maritime Self-Defence Force UH-60J (see also SH-60J/K/L) Royal Jordanian Air Force Latvian Air Force - UH-60M (2 received, another 2 on order) Lithuanian Air Force - UH-60M (4 on order; deliveries to begin in late 2024.) Malaysian Army - UH-60A+ (4 on lease, deliveries to begin in 2023) This lease was canceled in November 2024 due to budget issues. Royal Malaysian Air Force Mexican Air Force Mexican Navy Royal Moroccan Gendarmerie Philippine Air Force S-70i (21) (27 on order) Polish Special Forces - 4 S-70i helicopters (4 on order) Portuguese Air Force - UH-60A (9 ordered for aerial firefighting) Two received as of 2023. Royal Saudi Air Force Royal Saudi Land Forces Saudi Arabian National Guard Royal Saudi Navy Republic of Korea Air Force Republic of Korea Army UH-60P Republic of Korea Navy Slovak Air Force Swedish Air Force Republic of China Air Force Republic of China Army Republic of China Navy Royal Thai Army UH-60L; UH-60M Royal Thai Air Force Royal Thai Navy (see SH-60) Tunisian Air Force Turkish Air Force - (6 T-70s on order) First unit delivered in January 2023. Turkish Land Forces- 22+ T-70 ordered. First delivered Special Forces Command(Turkey) - 6 T-70s ordered with deliveries underway. Turkish Naval Forces (see Sikorsky SH-60 Seahawk) Gendarmerie General Command (Turkey) - 14 T-70s ordered with 3 delivered. General Directorate of Security(Turkey) - 20 T-70s ordered General Directorate of Forestry(Turkey) - 3 T-70s ordered with 2 delivered. United Arab Emirates Air Force Main Directorate of Intelligence (Ukraine) - 2 UH-60As One more being crowdfunded for GUR by Czech supporters under the "Gift for Putin" (Dárek pro-Putina) initiative. United States Air Force (see HH-60) United States Army - The U.S. Army has a stated requirement for 2,135 aircraft. United States Navy (see SH-60) United States Coast Guard (see MH-60) United States Department of State United States Department of Homeland Security Former operators - until August 2021 Afghan Air Force Royal Hong Kong Auxiliary Air Force Government Flying Service Accidents From 1981 to 1987, five Black Hawks crashed (killing or injuring all on board) while flying near radio broadcast towers because their electromagnetic emissions disrupted the helicopters' flight control systems. The Black Hawk helicopters were not hardened against high-intensity radiated fields, contrary to the SH-60 Seahawk Navy version. The pilots were instructed to fly away from emitters, and, in the long term, shielding was increased and backup systems were installed. On 29 July 1992, one Australian Army Black Hawk collided into terrain near Oakey Army Aviation Centre. Killing two occupants. On 3 March 1994, a UH-60 helicopter of the 15th Fighter Wing, Republic of Korea Air Force (ROKAF) exploded above Yongin, Gyeonggi-do, killing all of the six personnel on board, including General Cho Kun-hae, then Chief of the Air Staff of South Korea. On 14 April 1994, two U.S. Army UH-60 Black Hawks in northern Iraq were shot down in a friendly fire incident by U.S. Air Force F-15 fighter jets patrolling the northern no-fly zone that had been imposed after the 1991 Gulf War, in which all twenty-six crew and passengers were killed. The pilots of the U.S. Air Force F-15s misidentified the U.S. Army Black Hawk helicopters as enemy Mil Mi-24 "Hind" helicopters. On 12 June 1996, two Australian Army Black Hawks collided during an Army nighttime special forces counter-terrorism exercise resulting in the death of eighteen soldiers - fifteen members of the SASR and three from the 5th Aviation Regiment. On 12 February 2001, two Black Hawks from Schofield Barracks, Hawaii collided during NVG formation flight training, causing loss of both aircraft, six deaths and 11 injured soldiers. On 12 February 2004, one Australian Army Black Hawk collided into terrain in the vicinity of Mount Walker, Queensland following contact between the tail rotor and a tree. The airframe was written off however there were no deaths - six out of the eight occupants received injuries. On 26 September 2004, a U.S. Army Black Hawk crashed taking off from Tallil Airbase (Nasiriyah Airport), Iraq. The crew of four was rescued. On 29 November 2006, one Australian Army Black Hawk crashed into and subsequently slid off the deck of HMAS Kanimbla sinking into deep waters off the coast of Fiji whilst conducting a training flight. The sinking resulted in the deaths of two soldiers - one pilot from the 5th Aviation Regiment, and one trooper from the SASR. On 10 March 2015, a UH-60 from Eglin Air Force Base crashed off the coast of the Florida Panhandle near the base. All eleven on board were killed. On 16 February 2018, UH-60M helicopter deployed by the Mexican Air Force to Oaxaca after an earthquake, crashed into a group of people while attempting to land. On 2 January 2020, a UH-60M helicopter of the Republic of China Air Force (ROCAF) in Taiwan, crashed on a mountainside, killing eight people on board, including General Shen Yi-ming, chief of the general staff of Republic of China's armed forces. On 23 June 2021, a Philippine Air Force S-70i crashed in Capas town in Tarlac during a night flight training, killing all 6 crew members. The unit was newly delivered in November of the previous year or only almost 8 months old. On 22 February 2022, two Utah National Guard Black Hawk helicopters crashed at the Snowbird, Utah ski resort during a training exercise. One Black Hawk was overcome by whiteout conditions caused by the downdraft in the snow, and crashed, causing parts of the rotor blades to strike the other helicopter, forcing a hard landing. There were no major injuries to the crew or skiers. On July 16, 2022, one Mexican Navy Black Hawk crashed at Sinaloa, killing 14 marines on board. In September 2022, a Black Hawk operated by the Taliban crashed during a training exercise in Kabul, killing three. On February 15, 2023, a Black Hawk crashed killing two members of the Tennessee National Guard in Huntsville, Alabama. On 29 March 2023, two US Army Black Hawk medical helicopters crashed during a training mission over Kentucky. All nine soldiers aboard were killed. The cause of the crash is under investigation. On 10 November 2023, a US Army Black Hawk crashed off the coast of Cyprus in the Mediterranean Sea. All 5 soldiers aboard were killed. On 11 September 2024, an IDF Black Hawk crashed in Rafah, Gaza during a medical evacuation mission. Two soldiers were killed, and seven others were injured. Specifications (UH-60M)	Technology	Specific aircraft	null
37947	https://en.wikipedia.org/wiki/Violet%20%28color%29	Violet (color)	Violet is the color of light at the short wavelength end of the visible spectrum. It is one of the seven colors that Isaac Newton labeled when dividing the spectrum of visible light in 1672. Violet light has a wavelength between approximately 380 and 435 nanometers. The color's name is derived from the Viola genus of flowers. In the RGB color model used in computer and television screens, violet is produced by mixing red and blue light, with more blue than red. In the RYB color model historically used by painters, violet is created with a combination of red and blue pigments and is located between blue and purple on the color wheel. In the CMYK color model used in printing, violet is created with a combination of magenta and cyan pigments, with more magenta than cyan. On the RGB/CMY(K) color wheel, violet is located between blue and magenta. Violet is closely associated with purple. In optics, violet is a spectral color (referring to the color of different single wavelengths of light), whereas purple is the color of various combinations of red and blue (or violet) light, some of which humans perceive as similar to violet. In common usage, both terms are used to refer to a variety of colors between blue and red in hue. Violet has a long history of association with royalty, originally because Tyrian purple dye was extremely expensive in antiquity. The emperors of Rome wore purple togas, as did the Byzantine emperors. During the Middle Ages, violet was worn by bishops and university professors and was often used in art as the color of the robes of the Virgin Mary. In Chinese painting, the color violet represents the "unity transcending the duality of Yin and yang" and "the ultimate harmony of the universe". In New Age thinking, purple and/or violet is associated with the crown chakra. One European study suggests that violet is the color people most often associate with extravagance, individualism, vanity and ambiguity. Etymology and definitions The word violet as a color name derives from the Middle English and Old French violete, in turn from the Latin viola, the name of the violet flower. The first recorded use as a color name in English was in 1370. Relationship to purple Violet is closely associated with purple. In optics, violet is a spectral color: It refers to the color of any different single wavelength of light on the short wavelength end of the visible spectrum (between approximately 380 and 435 nanometers), whereas purple is the color of various combinations of red, blue and violet light, some of which some humans perceive as similar to violet. In common usage, both terms are used to refer to a variety of colors between blue and red in hue. Historically, violet has tended to be used for bluer hues and purple for redder hues. In the traditional color wheel used by painters, violet and purple are both placed between red and blue, with violet being closer to blue. In science Optics Violet is at one end of the spectrum of visible light, between blue light, which has a longer wavelength, and ultraviolet light, which has a shorter wavelength and is not visible to humans. Violet wavelengths are between approximately 380 and 435 nanometers. Violet objects often appear dark, because human vision becomes less sensitive at wavelengths this short. The reason why to (typical trichromat) humans violet light appears slightly reddish compared to spectral blue (despite spectral red being at the other end of the visible spectrum) is, according to the opponent process hypothesis of color vision, that the S-cone type (i.e. the one most sensitive to short wavelengths) contributes some red to the red-versus-green opponent channel (which at the longer blue wavelengths gets counteracted by the M-cone type). Computer and television screens, using the RGB color model, cannot produce actual violet light and instead mimic it by combining blue light at high intensity with red light at less intensity. Violet objects are normally composed-light violet. Objects reflecting spectral violet appear very dark, because human vision is relatively insensitive to those wavelengths. Monochromatic lamps emitting spectral-violet wavelengths can be roughly approximated by the color named electric violet, which is a composed-light violet producing a similar effect to the human eye. Chemistry – pigments and dyes The earliest violet pigments used by humans, found in prehistoric cave paintings, were made from the minerals manganese and hematite. Manganese is still used today by the Aranda people, a group of indigenous Australians, as a traditional pigment for coloring the skin during rituals. It is also used by the Hopi Indians of Arizona to color ritual objects. The most famous violet-purple dye in the ancient world was Tyrian purple, made from a type of sea snail called the murex, found around the Mediterranean. In western Polynesia, residents of the islands made a violet dye similar to Tyrian purple from the sea urchin. In Central America, the inhabitants made a dye from a different sea snail, the purpura, found on the coasts of Costa Rica and Nicaragua. The Mayans used this color to dye fabric for religious ceremonies, and the Aztecs used it for paintings of ideograms, where it symbolized royalty. During the Middle Ages, most artists made purple or violet on their paintings by combining red and blue pigments; usually blue azurite or lapis-lazuli with red ochre, cinnabar or minium. They also combined lake colors by mixing dye with powder; woad or indigo dye for blue and cochineal dye for red. Orcein, or purple moss, was another common violet dye. It was known to the ancient Greeks and Hebrews, was made from a Mediterranean lichen called archil or dyer's moss (Roccella tinctoria), combined with an ammoniac, usually urine. Orcein began to achieve popularity again in the 19th century, when violet and purple became the color of demi-mourning, worn after a widow or widower had worn black for a certain time, before he or she returned to wearing ordinary colors. In the 18th century, chemists in England, France and Germany began to create the first synthetic dyes. Two synthetic purple dyes were invented at about the same time. Cudbear is a dye extracted from orchil lichens that can be used to dye wool and silk, without the use of mordant. Cudbear was developed by Cuthbert Gordon of Scotland: production began in 1758, The lichen is first boiled in a solution of ammonium carbonate. The mixture is then cooled and ammonia is added and the mixture is kept damp for 3–4 weeks. Then the lichen is dried and ground to powder. The manufacture details were carefully protected, with a ten-foot high wall built around the manufacturing facility, and staff consisting of Highlanders sworn to secrecy. French purple was developed in France at about the same time. The lichen is extracted by urine or ammonia. Then the extract is acidified, the dissolved dye precipitates and is washed. Then it is dissolved in ammonia again, the solution is heated in air until it becomes purple, then it is precipitated with calcium chloride; the resulting dye was more solid and stable than other purples. Cobalt violet is a synthetic pigment that was invented in the second half of the 19th century, and is made by a similar process as cobalt blue, cerulean blue and cobalt green. It is the violet pigment most commonly used today by artists, along with manganese violet. Mauveine, also known as aniline purple and Perkin's mauve, was the first synthetic organic chemical dye, discovered serendipitously during an attempt to make quinine in 1856. Its chemical name is 3-amino-2,±9-dimethyl-5-phenyl-7-(p-tolylamino) phenazinium acetate. In the 1950s, a new family of violet synthetic organic pigments called quinacridones came onto the market. They had originally been discovered in 1896, synthesized in 1936 and manufactured in the 1950s. The colors in the group range from deep red to violet in color, and have the molecular formula C20H12N2O2. They have strong resistance to sunlight and washing, and are used in oil paints, watercolors and acrylics, as well as in automobile coatings and other industrial coatings. Zoology Botany In history and art Prehistory and antiquity Violet is one of the oldest colors used by humans. Traces of very dark violet, made by grinding the mineral manganese, mixed with water or animal fat and then brushed on the cave wall or applied with the fingers, are found in the prehistoric cave art in Pech Merle, in France, dating back about 25,000 years. It has also been found in the cave of Altamira and Lascaux. It was sometimes used as an alternative to black charcoal. Sticks of manganese, used for drawing, have been found at sites occupied by Neanderthals in France and Israel. From the grinding tools at various sites, it appears it may also have been used to color the body and to decorate animal skins. More recently, the earliest dates on cave paintings have been pushed back farther than 35,000 years. Hand paintings on rock walls in Australia may be even older, dating back as far as 50,000 years. Berries of the genus rubus, such as blackberries, were a common source of dyes in antiquity. The ancient Egyptians made a kind of violet dye by combining the juice of the mulberry with crushed green grapes. The Roman historian Pliny the Elder reported that the Gauls used a violet dye made from bilberry to color the clothing of slaves. These dyes made a satisfactory purple, but it faded quickly in sunlight and when washed. Middle Ages and Renaissance Violet and purple retained their status as the color of emperors and princes of the church throughout the long rule of the Byzantine Empire. While violet was worn less frequently by Medieval and Renaissance kings and princes, it was worn by the professors of many of Europe's new universities. Their robes were modeled after those of the clergy, and they often wore square violet caps and violet robes, or black robes with violet trim. Violet also played an important part in the religious paintings of the Renaissance. Angels and the Virgin Mary were often portrayed wearing violet robes. The 15th-century Florentine painter Cennino Cennini advised artists: "If you want to make a lovely violet colour, take fine lacca, ultramarine blue (the same amount of the one as of the other)..." For fresco painters, he advised a less-expensive version, made of a mixture of blue indigo and red hematite. 18th and 19th centuries In the 18th century, purple was a color worn by royalty, aristocrats and other wealthy people. Good-quality purple fabric was too expensive for ordinary people. The first cobalt violet, the intensely red-violet cobalt arsenate, was highly toxic. Although it persisted in some paint lines into the 20th century, it was displaced by less toxic cobalt compounds such as cobalt phosphate. Cobalt violet appeared in the second half of the 19th century, broadening the palette of artists with its range of purple colors. Cobalt violet was used by Paul Signac (1863–1935), Claude Monet (1840–1926) and Georges Seurat (1859–1891). Today, cobalt ammonium phosphate, cobalt lithium phosphate and cobalt phosphate are available for use by artists. Cobalt ammonium phosphate is the most reddish of the three. Cobalt phosphate is available in two varieties — a deep less saturated blueish type and a lighter and brighter somewhat more reddish type. Cobalt lithium phosphate is a saturated lighter-valued bluish violet. A color similar to cobalt ammonium phosphate, cobalt magnesium borate, was introduced in the later 20th century but was not deemed sufficiently lightfast for artistic use. Cobalt violet is the only truly lightfast purple pigment with relatively strong color saturation. All other light-stable purple pigments are dull by comparison. The high price of the pigment and the toxicity of cobalt have limited its use. In the 1860s, the popularity of using violet colors suddenly rose among painters and other artists. For example, Vincent van Gogh (1853–1890) was an avid student of color theory. He used violet in many of his paintings of the 1880s, including his paintings of irises and the swirling and mysterious skies of his starry night paintings, and often combined it with its complementary color, yellow. In his painting of his bedroom in Arles (1888), he used several sets of complementary colors; violet and yellow, red and green and orange and blue. In a letter about the painting to his brother Theo, he wrote, "The color here...should be suggestive of sleep and repose in general....The walls are a pale violet. The floor is of red tiles. The wood of the bed and the chairs are fresh butter yellow, the sheet and the pillows light lemon green. The bedspread bright scarlet. The window green. The bed table orange. The bowl blue. The doors lilac....The painting should rest the head or the imagination." In 1856, a young British chemist named William Henry Perkin was trying to make a synthetic quinine. His experiments produced instead an unexpected residue, which turned out to be the first synthetic aniline dye, a deep purple color called mauveine, or abbreviated simply to mauve (the dye being named after the lighter color of the mallow [mauve] flower). Used to dye clothes, it became extremely fashionable among the nobility and upper classes in Europe, particularly after Queen Victoria wore a silk gown dyed with mauveine to the Royal Exhibition of 1862. Prior to Perkin's discovery, mauve was a color which only the aristocracy and rich could afford to wear. Perkin developed an industrial process, built a factory, and produced the dye by the ton so almost anyone could wear mauve. It was the first of a series of modern industrial dyes which completely transformed both the chemical industry and fashion. 20th and 21st centuries Violet or purple neckties became popular at the end of the first decade of the 21st century, particularly among political and business leaders. In culture Popularity In a European survey, three percent of respondents said violet is their favorite color, ranking it behind blue, green, red, black and yellow (in that order), and tied with orange. Ten percent called it their least favorite color; brown, pink and gray were more unpopular. Royalty and luxury Because of its status as the color of Roman emperors, monarchs and princes, purple and violet are often associated with luxury. Certain luxury goods, such as watches and jewelry, are often placed in boxes lined with violet velvet, since violet is the complementary color of yellow and shows gold to best advantage. Vanity, extravagance and individualism While violet is the color of humility in the symbolism of the Catholic Church, it has exactly the opposite meaning in general society. A European poll in 2000 showed it was the color most commonly associated with vanity. As a color that rarely exists in nature and so attracts attention, it is seen as a color of individualism and extravagance. Heian period In Japan, violet was a popular color introduced into dress during the Heian period (794–1185). The dye was made from the root of the alkanet plant (Anchusa officinalis), known as murasaki in Japanese. At about the same time, Japanese painters began to use a pigment made from the same plant. New Age The "New Age Prophetess", Alice Bailey, in her system called the Seven Rays which classifies humans into seven different metaphysical psychological types, the "seventh ray" of "Ceremonial Order" is represented by the color violet. People who have this metaphysical psychological type are said to be "on the Violet Ray". In the Ascended Master Teachings, the color violet is used to represent the Ascended Master St. Germain. The Invocation of the Violet Flame is a system of meditation practice used in the "I AM" Activity and by the Church Universal and Triumphant (both Ascended Master Teaching religions). Religion In the Roman Catholic church, violet is worn by bishops and archbishops, red by cardinals and white by the Pope. Ordinary priests wear black. As in many other Western churches, violet is the liturgical color of Advent and Lent, which respectively celebrate the expectant waiting and preparation for the celebration of the Crucifixion of Jesus and the time for penance and/or mourning. A stained glass window installed in the early 1920s in the Cathedral of Our Lady of the Angels in Los Angeles depicts God the Father wearing a violet robe. After the Vatican II Council, which modified many of the rules of the Catholic church, priests began to wear violet robes when celebrating masses for the dead. Black was no longer used, since it was the color of mourning outside the church and deemed inappropriate in a religious ceremony. In Hinduism, violet is used to symbolically represent the seventh, crown chakra (Sahasrara). Politics In the early 20th century, violet, white and gold were the colors of the women's suffrage movement in the United States, seeking the right to vote for women. The colors were said to represent liberty and dignity. For this reason, the postage stamp issued in 1936 to honor Susan B. Anthony, a prominent leader of the suffrage movement in the United States, was colored the reddish tone of violet sometimes known as red-violet. In 1908, Emmeline Pethick-Lawrence, co-editor of the Women's Social and Political Union (WSPU) newspaper, designed the color scheme for the suffragette movement in Britain and Ireland, with violet for loyalty and dignity, white for purity and green for hope. The pan-European movement Volt Europa and its national subsidiary parties use violet in their uniforms. A small New Age political party in Germany with about 1,150 members is called The Violet Party. It believes in direct democracy, a guaranteed minimum income and politics based on spirituality. It was founded in Dortmund in 2001. Lesbianism Violet flowers and their color became symbolically associated with lesbianism. It was used as a special code by lesbians and bisexual women for self-identification and also to communicate support for the sexual preference. This connection originates from the poet Sappho and fragments of her poems. In one poem, she describes a lost partner wearing a garland of "violet tiaras, braided rosebuds, dill and crocus twined around" her neck. In another fragment, she recalls her partner as having "put around yourself [many wreaths] of violets and roses." The labrys lesbian flag, created in 1999 by graphic designer Sean Campbell, features a labrys superimposed on the inverted black triangle set against a violet background. Flags	Physical sciences	Colors	Physics
37948	https://en.wikipedia.org/wiki/Purple	Purple	Purple is a color similar in appearance to violet light. In the RYB color model historically used in the arts, purple is a secondary color created by combining red and blue pigments. In the CMYK color model used in modern printing, purple is made by combining magenta pigment with either cyan pigment, black pigment, or both. In the RGB color model used in computer and television screens, purple is created by mixing red and blue light in order to create colors that appear similar to violet light. Purple has long been associated with royalty, originally because Tyrian purple dye—made from the secretions of sea snails—was extremely expensive in antiquity. Purple was the color worn by Roman magistrates; it became the imperial color worn by the rulers of the Byzantine Empire and the Holy Roman Empire, and later by Roman Catholic bishops. Similarly in Japan, the color is traditionally associated with the emperor and aristocracy. According to contemporary surveys in Europe and the United States, purple is the color most often associated with rarity, royalty, luxury, ambition, magic, mystery, piety and spirituality. When combined with pink, it is associated with eroticism, femininity, and seduction. Etymology and definitions The modern English word purple comes from the Old English purpul, which derives from Latin purpura, which, in turn, derives from the Greek (porphura), the name of the Tyrian purple dye manufactured in classical antiquity from a mucus secreted by the spiny dye-murex snail. The first recorded use of the word purple dates to the late 900s AD. In art, history, and fashion In prehistory and the ancient world Purple first appeared in prehistoric art during the Neolithic era. The artists of Pech Merle cave and other Neolithic sites in France used sticks of manganese and hematite powder to draw and paint animals and the outlines of their own hands on the walls of their caves. These works have been dated to between 16,000 and 25,000 BC. Purple textiles, dating back to the early second millennium BCE, were found in Syria, making them the oldest known purple textiles in the world. These findings include textiles from a burial site in Chagar Bazar, dating back to the 18th-16th centuries BCE, as well as preserved textile samples discovered in gypsum at the Royal Palace of Qatna. As early as the 15th century BC, the citizens of Sidon and Tyre, two cities on the coast of Ancient Phoenicia (present day Lebanon), were producing purple dye from a sea snail called the spiny dye-murex. Clothing colored with the Tyrian dye was mentioned in both the Iliad of Homer and the Aeneid of Virgil. The deep, rich purple dye made from this snail became known as Tyrian purple. The process of making the dye was long, difficult and expensive. Thousands of the tiny snails had to be found, their shells cracked, the snail removed. Mountains of empty shells have been found at the ancient sites of Sidon and Tyre. The snails were left to soak, then a tiny gland was removed and the juice extracted and put in a basin, which was placed in the sunlight. There, a remarkable transformation took place. In the sunlight the juice turned white, then yellow-green, then green, then violet, then a red which turned darker and darker. The process had to be stopped at exactly the right time to obtain the desired color, which could range from a bright crimson to a dark purple, the color of dried blood. Then either wool, linen or silk would be dyed. The exact hue varied between crimson and violet, but it was always rich, bright and lasting. Tyrian purple became the color of kings, nobles, priests and magistrates all around the Mediterranean. It was mentioned in the Hebrew Bible (Old Testament); in the Book of Exodus, God instructs Moses to have the Israelites bring him an offering including cloth "of blue, and purple, and scarlet," to be used in the curtains of the Tabernacle and the garments of priests. The term used for purple in the 4th-century Latin Vulgate version of the Bible passage is purpura or Tyrian purple. In the Iliad of Homer, the belt of Ajax is purple, and the tails of the horses of Trojan warriors are dipped in purple. In the Odyssey, the blankets on the wedding bed of Odysseus are purple. In the poems of Sappho (6th century BC) she celebrates the skill of the dyers of the Greek kingdom of Lydia who made purple footwear, and in the play of Aeschylus (525–456 BC), Queen Clytemnestra welcomes back her husband Agamemnon by decorating the palace with purple carpets. In 950 BC, King Solomon was reported to have brought artisans from Tyre to provide purple fabrics to decorate the Temple of Jerusalem. Alexander the Great (when giving imperial audiences as the basileus of the Macedonian Empire), the basileus of the Seleucid Empire, and the kings of Ptolemaic Egypt all wore Tyrian purple. The Roman custom of wearing purple togas may have come from the Etruscans; an Etruscan tomb painting from the 4th century BC shows a nobleman wearing a deep purple and embroidered toga. In Ancient Rome, the Toga praetexta was an ordinary white toga with a broad purple stripe on its border. It was worn by freeborn Roman boys who had not yet come of age, curule magistrates, certain categories of priests, and a few other categories of citizens. The Toga picta was solid purple, embroidered with gold. During the Roman Republic, it was worn by generals in their triumphs, and by the Praetor Urbanus when he rode in the chariot of the gods into the circus at the Ludi Apollinares. During the Empire, the toga picta was worn by magistrates giving public gladiatorial games, and by the consuls, as well as by the emperor on special occasions. During the Roman Republic, when a triumph was held, the general being honored wore an entirely purple toga bordered in gold, and Roman Senators wore a toga with a purple stripe. However, during the Roman Empire, purple was more and more associated exclusively with the emperors and their officers. Suetonius claims that the early emperor Caligula had the King of Mauretania murdered for the splendour of his purple cloak, and that Nero forbade the use of certain purple dyes. In the late empire the sale of purple cloth became a state monopoly protected by the death penalty. According to the New Testament, Jesus Christ, in the hours leading up to his crucifixion, was dressed in purple (πορφύρα: porphura) by the Roman garrison to mock his claim to be 'King of the Jews'. The actual color of Tyrian purple seems to have varied from a reddish to a bluish purple. According to the Roman writer Vitruvius, (1st century BC), the murex shells coming from northern waters, probably Bolinus brandaris, produced a more bluish color than those of the south, probably Hexaplex trunculus. The most valued shades were said to be those closer to the color of dried blood, as seen in the mosaics of the robes of the Emperor Justinian in Ravenna. The chemical composition of the dye from the murex is close to that of the dye from indigo, and indigo was sometimes used to make a counterfeit Tyrian purple, a crime which was severely punished. What seems to have mattered about Tyrian purple was not its color, but its luster, richness, its resistance to weather and light, and its high price. In modern times, Tyrian purple has been recreated, at great expense. When the German chemist Paul Friedander tried to recreate Tyrian purple in 2008, he needed twelve thousand mollusks to create 1.4 ounces of dye, enough to color a handkerchief. In the year 2000, a gram of Tyrian purple made from ten thousand mollusks according to the original formula cost two thousand euros. China In ancient China, purple was obtained not through the Mediterranean mollusc, but purple gromwell. The dye obtained did not easily adhere to fabrics, making purple fabrics expensive. Purple became a fashionable color in the state of Qi (齊, 1046 BC–221 BC) because its ruler, Duke Huan of Qi, developed a preference for it. As a result, the price of purple fabric was over five times that of plain fabric. His minister, Guan Zhong (管仲), eventually convinced him to relinquish this preference. China was the first culture to develop a synthetic purple color. An old hypothesis suggested links between the Chinese purple and blue and Egyptian blue, however, molecular structure analysis and evidence such as the absence of lead in Egyptian blue and the lack of examples of Egyptian blue in China, argued against the hypothesis. The use of quartz, barium, and lead components in ancient Chinese glass and Han purple and Han blue has been used to suggest a connection between glassmaking and the manufacture of pigments, and to prove the independence of the Chinese invention. Taoist alchemists may have developed Han purple from their knowledge of glassmaking. Lead is used by the pigment maker to lower the melting point of the barium in Han Purple. Purple was regarded as a secondary color in ancient China. In classical times, secondary colors were not as highly prized as the five primary colors of the Chinese spectrum, and purple was used to allude to impropriety, in contrast to crimson, which was deemed a primary color and symbolized legitimacy. Nevertheless, by the 6th century AD, purple was ranked above crimson. Several changes to the ranks of colors occurred after that time. Purple in the Byzantine Empire and Carolingian Europe Through the early Christian era, the rulers of the Byzantine Empire continued the use of purple as the imperial color, for diplomatic gifts, and even for imperial documents and the pages of the Bible. Gospel manuscripts were written in gold lettering on parchment that was colored Tyrian purple. Empresses gave birth in the Purple Chamber, and the emperors born there were known as "born to the purple," to separate them from emperors who won or seized the title through political intrigue or military force. Bishops of the Byzantine church wore white robes with stripes of purple, while government officials wore squares of purple fabric to show their rank. In western Europe, the Emperor Charlemagne was crowned in 800 wearing a mantle of Tyrian purple, and was buried in 814 in a shroud of the same color, which still exists (see below). However, after the fall of Constantinople to the Ottoman Turks in 1453, the color lost its imperial status. The great dye works of Constantinople were destroyed, and gradually scarlet, made with dye from the cochineal insect, became the royal color in Europe. The Middle Ages and Renaissance In 1464, Pope Paul II decreed that cardinals should no longer wear Tyrian purple, and instead wear scarlet, from kermes and alum, since the dye from Byzantium was no longer available. Bishops and archbishops, of a lower status than cardinals, were assigned the color purple, but not the rich Tyrian purple. They wore cloth dyed first with the less expensive indigo blue, then overlaid with red made from kermes dye. While purple was worn less frequently by Medieval and Renaissance kings and princes, it was worn by the professors of many of Europe's new universities. Their robes were modeled after those of the clergy, and they often wore square/violet or purple/violet caps and robes, or black robes with purple/violet trim. Purple/violet robes were particularly worn by students of divinity. Purple and violet also played an important part in the religious paintings of the Renaissance. Angels and the Virgin Mary were often portrayed wearing purple or violet robes. 18th and 19th centuries In the 18th century, purple was still worn on occasion by Catherine the Great and other rulers, by bishops and, in lighter shades, by members of the aristocracy, but rarely by ordinary people, because of its high cost. But in the 19th century, that changed. In 1856, an eighteen-year-old British chemistry student named William Henry Perkin was trying to make a synthetic quinine. His experiments produced instead the first synthetic aniline dye, a purple shade called mauveine, shortened simply to mauve. It took its name from the mallow flower, which is the same color. The new color quickly became fashionable, particularly after Queen Victoria wore a silk gown dyed with mauveine to the Royal Exhibition of 1862. Prior to Perkin's discovery, mauve was a color which only the aristocracy and rich could afford to wear. Perkin developed an industrial process, built a factory, and produced the dye by the ton, so almost anyone could wear mauve. It was the first of a series of modern industrial dyes which completely transformed both the chemical industry and fashion. Purple was popular with the pre-Raphaelite painters in Britain, including Arthur Hughes, who loved bright colors and romantic scenes. 20th and 21st centuries At the turn of the century, purple was a favorite color of the Austrian painter Gustav Klimt, who flooded his pictures with sensual purples and violets. In the 20th century, purple retained its historic connection with royalty; George VI (1896–1952), wore purple in his official portrait, and it was prominent in every feature of the coronation of Elizabeth II in 1953, from the invitations to the stage design inside Westminster Abbey. But at the same time, it was becoming associated with social change; with the Women's Suffrage movement for the right to vote for women in the early decades of the century, with Feminism in the 1970s, and with the psychedelic drug culture of the 1960s. In the early 20th century, purple, green, and white were the colors of the Women's Suffrage movement, which fought to win the right to vote for women, finally succeeding with the 19th Amendment to the U.S. Constitution in 1920. Later, in the 1970s, in a tribute to the Suffragettes, it became the color of the women's liberation movement. In the concentration camps of Nazi Germany, prisoners who were members of non-conformist religious groups, such as the Jehovah's Witnesses, were required to wear a purple triangle. During the 1960s and early 1970s, it was also associated with counterculture, psychedelics, and musicians like Jimi Hendrix with his 1967 song "Purple Haze", or the English rock band of Deep Purple which formed in 1968. Later, in the 1980s, it was featured in the song and album Purple Rain (1984) by the American musician Prince. The Purple Rain Protest was a protest against apartheid that took place in Cape Town, South Africa on 2 September 1989, in which a police water cannon with purple dye sprayed thousands of demonstrators. This led to the slogan The Purple Shall Govern. The violet or purple necktie became very popular at the end of the first decade of the 21st century, particularly among political and business leaders. It combined the assertiveness and confidence of a red necktie with the sense of peace and cooperation of a blue necktie, and it went well with the blue business suit worn by most national and corporate leaders. In science and nature Optics The meanings of the color terms violet and purple varies even among native speakers of English, for example between United Kingdom and United States. Optics research on purple and violet contains contributions of authors from different countries and different native languages, it is likely to be inconsistent in the use and meaning of the two colors. According to some speakers/authors of English, purple, unlike violet, is not one of the colors of the visible spectrum. It was not one of the colors of the rainbow identified by Isaac Newton. According to some authors, purple does not have its own wavelength of light. For this reason, it is sometimes called a non-spectral color. It exists in culture and art, but not, in the same way that violet does, in optics. According to some speakers of English, purple is simply a combination, in various proportions, of two primary colors, red and blue. According to other speakers of English, the same range of colors is called violet. In some textbooks of color theory, and depending on the geographical-cultural origin of the author, a "purple" is defined as any non-spectral color between violet and red (excluding violet and red themselves). In that case, the spectral colors violet and indigo would not be shades of purple. For other speakers of English, these colors are shades of purple. In the traditional color wheel long used by painters, purple is placed between crimson and violet. However, also here there is much variation in color terminology depending on cultural background of the painters and authors, and sometimes the term violet is used and placed in between red and blue on the traditional color wheel. In a slightly different variation, on the color wheel, purple is placed between magenta and violet. This shade is sometimes called electric purple (see shades of purple). In the RGB color model, named for the colors red, green, and blue, used to create all the colors on a computer screen or television, the range of purples is created by mixing red and blue light of different intensities on a black screen. The standard HTML color purple is created by red and blue light of equal intensity, at a brightness that is halfway between full power and darkness. In color printing, purple is sometimes represented by the color magenta, or sometimes by mixing magenta with red or blue. It can also be created by mixing just red and blue alone, but in that case the purple is less bright, with lower saturation or intensity. A less bright purple can also be created with light or paint by adding a certain quantity of the third primary color (green for light or yellow for pigment). Relationship with violet Purple is closely associated with violet. In common usage, both refer to a variety of colors between blue and red in hue. Historically, purple has tended to be used for redder hues and violet for bluer hues. In optics, violet is a spectral color; it refers to the color of any different single wavelength of light on the short wavelength end of the visible spectrum, between approximately 380 and 450 nanometers, whereas purple is the color of various combinations of red, blue, and violet light, some of which humans perceive as similar to violet. On a chromaticity diagram, the straight line connecting the extreme spectral colors (red and violet) is known as the line of purples (or 'purple boundary'); it represents one limit of human color perception. The color magenta used in the CMYK printing process is near the center of the line of purples, but most people associate the term "purple" with a somewhat bluer tone, such as is displayed by the color "electric purple" (a color also directly on the line of purples), shown below. On the CIE xy chromaticity diagram, violet is on the curved edge in the lower left, while purples are on the straight line connecting the extreme colors red and violet; this line is known as the line of purples, or the purple line. Pigments Hematite and manganese are the oldest pigments used for the color purple. They were used by Neolithic artists in the form of sticks, like charcoal, or ground and powdered and mixed with fat, and used as a paint. Hematite is a reddish iron oxide which, when ground coarsely, makes a purple pigment. One such pigment is caput mortuum, whose name is also used in reference to mummy brown. The latter is another pigment containing hematite and historically produced with the use of mummified corpses. Some of its compositions produce a purple color and may be called "mummy violet". Manganese was also used in Roman times to color glass purple. Han purple was the first synthetic purple pigment, invented in China in about 700 BC. It was used in wall paintings and pottery and other applications. In color, it was very close to indigo, which had a similar chemical structure. Han purple was very unstable, and sometimes was the result of the chemical breakdown of Han blue. During the Middle Ages, artists usually made purple by combining red and blue pigments; most often blue azurite or lapis-lazuli with red ochre, cinnabar, or minium. They also combined lake colors made by mixing dye with powder; using woad or indigo dye for the blue, and dye made from cochineal for the red. Cobalt violet was the first modern synthetic color in the purple family, manufactured in 1859. It was found, along with cobalt blue, in the palette of Claude Monet, Paul Signac, and Georges Seurat. It was stable, but had low tinting power and was expensive, so quickly went out of use. Manganese violet was a stronger color than cobalt violet, and replaced it on the market. Quinacridone violet, one of a modern synthetic organic family of colors, was discovered in 1896 but not marketed until 1955. It is sold today under a number of brand names. Dyes The most famous purple dye in the ancient world was Tyrian purple, made from a type of sea snail called the murex, found around the Mediterranean. (See history section above). In western Polynesia, residents of the islands made a purple dye similar to Tyrian purple from the sea urchin. In Central America, the inhabitants made a dye from a different sea snail, the purpura, found on the coasts of Costa Rica and Nicaragua. The Mayans used this color to dye fabric for religious ceremonies, while the Aztecs used it for paintings of ideograms, where it symbolized royalty. In the Middle Ages, those who worked with blue and black dyes belonged to separate guilds from those who worked with red and yellow dyes, and were often forbidden to dye any other colors than those of their own guild. Most purple fabric was made by the dyers who worked with red, and who used dye from madder or cochineal, so Medieval violet colors were inclined toward red. Orcein, or purple moss, was another common purple dye. It was known to the ancient Greeks and Hebrews, and was made from a Mediterranean lichen called archil or dyer's moss (Roccella tinctoria), combined with an ammoniac, usually urine. Orcein began to achieve popularity again in the 19th century, when violet and purple became the color of demi-mourning, worn after a widow or widower had worn black for a certain time, before he or she returned to wearing ordinary colors. From the Middle Ages onward, purple dyes for the clothing of common people were often made from the blackberry or other red fruit of the genus rubus, or from the mulberry. All of these dyes were more reddish than bluish, and faded easily with washing and exposure to sunlight. A popular new dye which arrived in Europe from the New World during the Renaissance was made from the wood of the logwood tree (Haematoxylum campechianum), which grew in Spanish Mexico. Depending on the different minerals added to the dye, it produced a blue, red, black or, with the addition of alum, a purple color, it made a good color, but, like earlier dyes, it did not resist sunlight or washing. In the 18th century, chemists in England, France and Germany began to create the first synthetic dyes. Two synthetic purple dyes were invented at about the same time. Cudbear is a dye extracted from orchil lichens that can be used to dye wool and silk, without the use of mordant. Cudbear was developed by Dr Cuthbert Gordon of Scotland: production began in 1758, The lichen is first boiled in a solution of ammonium carbonate. The mixture is then cooled and ammonia is added and the mixture is kept damp for 3–4 weeks. Then the lichen is dried and ground to powder. The manufacture details were carefully protected, with a ten-feet high wall being built around the manufacturing facility, and staff consisting of Highlanders sworn to secrecy. French purple was developed in France at about the same time. The lichen is extracted by urine or ammonia. Then the extract is acidified, the dissolved dye precipitates and is washed. Then it is dissolved in ammonia again, the solution is heated in air until it becomes purple, then it is precipitated with calcium chloride; the resulting dye was more solid and stable than other purples. Cobalt violet is a synthetic pigment that was invented in the second half of the 19th century, and is made by a similar process as cobalt blue, cerulean blue and cobalt green. It is the violet pigment most commonly used today by artists. In spite of its name, this pigment produces a purple rather than violet color Mauveine, also known as aniline purple and Perkin's mauve, was the first synthetic organic chemical dye, discovered serendipitously in 1856. Its chemical name is 3-amino-2,±9-dimethyl-5-phenyl-7-(p-tolylamino)phenazinium acetate. Fuchsine was another synthetic dye made shortly after mauveine. It produced a brilliant fuchsia color. In the 1950s, a new family of purple and violet synthetic organic pigments called quinacridone came onto the market. It had originally been discovered in 1896, but were not synthesized until 1936, and not manufactured until the 1950s. The colors in the group range from deep red to bluish purple in color, and have the molecular formula C20H12N2O2. They have strong resistance to sunlight and washing, and are widely used today in oil paints, water colors, and acrylics, as well as in automobile coatings and other industrial coatings. Animals Anthocyanins Certain grapes, eggplants, pansies and other fruits, vegetables and flowers may appear purple due to the presence of natural pigments called anthocyanins. These pigments are found in the leaves, roots, stems, vegetables, fruits and flowers of all plants. They aid photosynthesis by blocking harmful wavelengths of light that would damage the leaves. In flowers, the purple anthocyanins help attract insects who pollinate the flowers. Not all anthocyanins are purple; they vary in color from red to purple to blue, green, or yellow, depending upon the level of their pH. Plants and flowers Purple needlegrass is the state grass of California. Microbiology Purple bacteria are bacteria that are phototrophic, that is, capable of producing energy through photosynthesis. In April 2007, it was suggested that early archaea may have used retinal, a purple pigment, instead of chlorophyll, to extract energy from the sun. If so, large areas of the ocean and shoreline would have been colored purple; this is called the Purple Earth hypothesis. Astronomy One of the stars in the Pleiades, called Pleione, is sometimes called Purple Pleione because, being a fast spinning star, it has a purple hue caused by its blue-white color being obscured by a spinning ring of electrically excited red hydrogen gas. The Purple Forbidden enclosure is a name used in traditional Chinese astronomy for those Chinese constellations that surround the north celestial pole. Geography Purple Mountain is located on the eastern side of Nanjing. Its peaks are often found enveloped in purple clouds at dawn and dusk, hence comes its name "Purple Mountain". The Purple Mountain Observatory is located there. Purple Mountain in County Kerry, Ireland, takes its name from the color of the shivered slate on its summit. Purple Mountain in Wyoming (el. ) is a mountain peak in the southern section of the Gallatin Range in Yellowstone National Park. Purple Mountain, Alaska Purple Mountain, Oregon Purple Mountain, Washington Purple Peak, Colorado Purple mountains phenomenon It has been observed that the greater the distance between a viewers eyes and mountains, the lighter and more blue or purple they will appear. This phenomenon, long recognized by Leonardo da Vinci and other painters, is called aerial perspective or atmospheric perspective. The more distant the mountains are, the less contrast the eye sees between the mountains and the sky. The bluish color is caused by an optical effect called Rayleigh scattering. The sunlit sky is blue because air scatters short-wavelength light more than longer wavelengths. Since blue light is at the short wavelength end of the visible spectrum, it is more strongly scattered in the atmosphere than long wavelength red light. The result is that the human eye perceives blue when looking toward parts of the sky other than the sun. At sunrise and sunset, the light is passing through the atmosphere at a lower angle, and traveling a greater distance through a larger volume of air. Much of the green and blue is scattered away, and more red light comes to the eye, creating the colors of the sunrise and sunset and making the mountains look purple. The phenomenon is referenced in the song "America the Beautiful", where the lyrics refer to "purple mountains' majesty" among other features of the United States landscape. A Crayola crayon called Purple Mountain Majesty in reference to the lyric was first formulated in 1993. Mythology Julius Pollux, a Greek grammarian who lived in the second century AD, attributed the discovery of purple to the Phoenician god and guardian of the city of Tyre, Heracles. According to his account, while walking along the shore with the nymph Tyrus, the god's dog bit into a murex shell, causing his mouth to turn purple. The nymph subsequently requested that Heracles create a garment for her of that same color, with Heracles obliging her demands giving birth to Tyrian purple. Associations and symbolism Royalty In Europe, since some Roman emperors wore a Tyrian purple (purpura) toga praetexta, purple has been the color most associated with power and royalty. The British Royal Family and other European royalty still use it as a ceremonial color on special occasions. In Japan, purple is associated with the emperor and Japanese aristocracy. Piety, faith, penitence, and theology In the West, purple or violet is a color often associated with piety and religious faith. In AD 1464, shortly after the Muslim conquest of Constantinople, which terminated the supply of Tyrian purple to Roman Catholic Europe, Pope Paul II decreed that cardinals should henceforth wear scarlet instead of purple, the scarlet being dyed with expensive cochineal. Bishops were assigned the color amaranth, being a pale and pinkish purple made then from a less-expensive mixture of indigo and cochineal. In the Latin liturgical rites of the Catholic liturgy, purple represents penitence; Anglican and Catholic priests wear a purple stole when they hear confession and a purple stole and chasuble during Advent and Lent. Since the Second Vatican Council of 1962–5, priests may wear purple vestments, but may still wear black ones, when officiating at funerals. The Roman Missal permits black, purple (violet), or white vestments for the funeral Mass. White is worn when a child dies before the age of reason. Students and faculty of theology also wear purple academic dress for graduations and other university ceremonies. Purple is also often worn by senior pastors of Protestant churches and bishops of the Anglican Communion. The color purple is also associated with royalty in Christianity, being one of the three traditional offices of Jesus Christ, i. e. king, although such a symbolism was assumed from the earlier Roman association or at least also employed by the ancient Romans. Vanity, extravagance, individualism In Europe and America, purple is the color most associated with vanity, extravagance, and individualism. Among the seven deadly sins, it represents pride. It is a color which is used to attract attention. The artificial, materialism and beauty Purple is the color most often associated with the artificial and the unconventional. It is the major color that occurs the least frequently in nature, and was the first color to be synthesized. Ambiguity and ambivalence Purple is the color most associated with ambiguity. Like other colors made by combining two primary colors, it is seen as uncertain and equivocal. Mourning In Britain, purple is sometimes associated with mourning. In Victorian times, close relatives wore black for the first year following a death ("deep mourning"), and then replaced it with purple or dark green trimmed with black. This is rarely practised today. In culture and society Cultures of Asian countries The Chinese word for purple, zi, is connected with the North Star, Polaris, or zi Wei in Chinese. In Chinese astrology, the North Star was the home of the Celestial Emperor, the ruler of the heavens. The area around the North Star is called the Purple Forbidden Enclosure in Chinese astronomy. For that reason the Forbidden City in Beijing was also known as the Purple Forbidden City (zi Jin cheng). Purple often represents "the highest," holiest, and "most sacred values" in China. In Taoism, purple is a transitional color and metaphysically between yin and yang. Purple was a popular color introduced into Japanese dress during the Heian period (794–1185). The dye was made from the root of the alkanet plant (Anchusa officinalis), also known as murasaki in Japanese. At about the same time, Japanese painters began to use a pigment made from the same plant. In Thailand, widows in mourning wear the color purple. Purple is also associated with Saturday on the Thai solar calendar. Cultures of Europe Ancient Rome Purple represented the height of Roman virtue and cultural values. Medieval Europe In medieval Europe, purple represented leadership and the king. In European alchemy during this time, "the 'precious purple tincture'" was a term for various substances alchemists hoped to create. The term and goal of the alchemists evoked kingliness, since the divine right of kings was also thought to aid the alchemists' future. Engineering The color purple plays a significant role in the traditions of engineering schools across Canada. Purple is also the color of the Engineering Corp in the British Military. Idioms and expressions Purple prose refers to pretentious or overly embellished writing. For example, a paragraph containing an excessive number of long and unusual words is called a purple passage. Born to the purple means someone who is born into a life of wealth and privilege. It originally was used to describe the rulers of the Byzantine Empire. A purple patch is a period of exceptional success or good luck. The origins are obscure, but it may refer to the symbol of success of the Byzantine Court. Bishops in Byzantium wore a purple patch on their costume as a symbol of rank. Purple haze refers to a state of mind induced by psychedelic drugs, particularly LSD. Wearing purple is a military slang expression in the U.S., Canada and the U.K. for an officer who is serving in a joint assignment with another service, such as an Army officer on assignment to the Navy. The officer is symbolically putting aside his or her traditional uniform color and exclusive loyalty to their service during the joint assignment, though in fact they continue to wear their own service's uniform. Purple squirrel is a term used by employment recruiters to describe a job candidate with precisely the right education, experience, and qualifications that perfectly fits a job's multifaceted requirements. The assumption is that the perfect candidate is as rare as a real-life purple squirrel. Military The Purple Heart is a United States military decoration awarded in the name of the President to those who have been wounded or killed during their service. Politics In United States politics, a purple state (typically a swing state) is a state roughly balanced between Republicans (generally symbolized by red in the 21st century) and Democrats (symbolized by blue). In the politics of the Netherlands, Purple () means a coalition government consisting of liberals and social democrats (symbolized by the colors blue and red, respectively), as opposed to the more common coalitions of the Christian Democrats with one of the other two. Between 1994 and 2002 there were two Purple cabinets, both led by Prime Minister Wim Kok. In the politics of Belgium, as with the Netherlands, a purple government includes liberal and social-democratic parties in coalition. Belgium was governed by Purple governments from 1999 to 2007 under the leadership of Prime Minister Guy Verhofstadt. Purple is the primary color used by many European and American political parties, including Volt Europa, the UK Independence Party, the Social Democrats in the Republic of Ireland, the Liberal People's Party in Norway, and the United States Pirate Party. The Left party in Germany, whose primary color is red, is traditionally portrayed in purple on election maps to distinguish it from the Social Democratic Party of Germany. In the United Kingdom, the color scheme for the suffragette movement in Britain and Ireland was designed with purple for loyalty and dignity, white for purity, and green for hope. Rhyme In the English language, the word "purple" has only one perfect rhyme, curple. Others are obscure perfect rhymes, such as hirple. Robert Burns rhymes purple with curple in his Epistle to Mrs. Scott. Examples of imperfect rhymes or non-word rhymes with purple: In the song Grace Kelly by Mika the word purple is rhymed with "hurtful". In his hit song "Dang Me", Roger Miller sings these lines: Sexuality Purple is sometimes associated with the lesbian, gay, bisexual, and transgender (LGBT) community. It is the symbolic color worn on Spirit Day, a commemoration that began in 2010 to show support for young people who are bullied because of their sexual orientation. Purple is closely associated with bisexuality, largely in part to the bisexual pride flag which combines pink – representing homosexuality – and blue – representing heterosexuality – to create the bisexual purple. The purple hand is another symbol sometimes used by the LGBT community during parades and demonstrations. Sports and games In Motorsport, purple is used to indicate the fastest times of the race. The National Basketball Association's Los Angeles Lakers, Phoenix Suns and Sacramento Kings use purple as their primary color. In the Indian Premier League, purple is the primary color of the Kolkata Knight Riders. In Major League Baseball, purple is one of the primary colors for the Colorado Rockies. In the National Football League, the Minnesota Vikings and Baltimore Ravens use purple as main colors. The Australian Football League's Fremantle Football Club use purple as one of their primary colors. In association football (soccer), Italian Serie A club ACF Fiorentina, Belgian Pro League club and former Europa League winner R.S.C. Anderlecht, French Ligue 1 club Toulouse FC and Ligue 2 club FC Istres, Spanish La Liga club Real Valladolid, Austrian Football Bundesliga club FK Austria Wien, Hungarian Nemzeti Bajnokság I club Újpest FC, Slovenian PrvaLiga club NK Maribor, former Romanian Liga I clubs FC Politehnica Timișoara and FC Argeș Pitești, Andorran Primera Divisió club CE Principat, German club Tennis Borussia Berlin, Italian club A.S.D. Legnano Calcio 1913, Swedish club Fässbergs IF, Japanese club Kyoto Sanga, Australian A-League Club Perth Glory and American Major League Soccer club Orlando City use purple as one of their primary colors. The Melbourne Storm from Australia's National Rugby League use purple as one of their primary colors. Costa Rica's Primera División soccer team Deportivo Saprissa's main color is purple (actually a burgundy like shade), and their nickname is the "Monstruo Morado", or "Purple Monster". In tennis, the official colors of the Wimbledon Championships are deep green and purple (traditionally called mauve). In American college athletics, Louisiana State University, Kansas State University, Texas Christian University, the University of Central Arkansas, Northwestern University, the University of Washington, and East Carolina University all have purple as one of their main team colors. The University of Western Ontario in London, Canada, and Bishop's University in Sherbrooke, Canada, have purple as one of its main team colors. Purple is the color of the ball in Snooker Plus with a 10-point value. In the game of pool, purple is the color of the 4-solid and the 12-striped balls. Business The British chocolate company Cadbury chose purple as it was Queen Victoria's favourite color. The company trademarked the color purple for chocolates with registrations in 1995 and 2004. However, the validity of these trademarks is the matter of an ongoing legal dispute following objections by Nestlé. In flags Purple or violet appear in the flags of only two modern sovereign nations, and are merely ancillary colors in both cases. The Flag of Dominica features a sisserou parrot, a national symbol, while the Flag of Nicaragua displays a rainbow in the center, as part of the coat of arms of Nicaragua. The lower band of the flag of the second Spanish republic (1931–39) was colored a tone of purple, to represent the common people as opposed to the red of the Spanish monarchy, unlike other nations of Europe where purple represented royalty and red represented the common people. In Japan, the prefecture of Tokyo's flag is purple, as is the flag of Ichikawa and other Japanese municipalities. Porpora, or purpure, a shade of purple, was added late to the list of colors of European heraldry. A purple lion was the symbol of the old Spanish Kingdom of León (910–1230), and it later appeared on the flag of Spain, when the Kingdom of Castile and Kingdom of León merged.	Physical sciences	Color terms	null
37955	https://en.wikipedia.org/wiki/Feldspar	Feldspar	Feldspar ( ; sometimes spelled felspar) is a group of rock-forming aluminium tectosilicate minerals, also containing other cations such as sodium, calcium, potassium, or barium. The most common members of the feldspar group are the plagioclase (sodium-calcium) feldspars and the alkali (potassium-sodium) feldspars. Feldspars make up about 60% of the Earth's crust and 41% of the Earth's continental crust by weight. Feldspars crystallize from magma as both intrusive and extrusive igneous rocks and are also present in many types of metamorphic rock. Rock formed almost entirely of calcic plagioclase feldspar is known as anorthosite. Feldspars are also found in many types of sedimentary rocks. Etymology The name feldspar derives from the German , a compound of the words ("field") and ("flake"). had long been used as the word for "a rock easily cleaved into flakes"; was introduced in the 18th century as a more specific term, referring perhaps to its common occurrence in rocks found in fields (Urban Brückmann, 1783) or to its occurrence as "fields" within granite and other minerals (René-Just Haüy, 1804). The change from to -spar was influenced by the English word spar, meaning a non-opaque mineral with good cleavage. Feldspathic refers to materials that contain feldspar. The alternate spelling, felspar, has fallen out of use. The term 'felsic', meaning light coloured minerals such as quartz and feldspars, is an acronymic word derived from feldspar and silica, unrelated to the obsolete spelling 'felspar'. Compositions The feldspar group of minerals consists of tectosilicates, silicate minerals in which silicon ions are linked by shared oxygen ions to form a three-dimensional network. Compositions of major elements in common feldspars can be expressed in terms of three endmembers: potassium feldspar (K-spar) endmember KAlSiO albite endmember NaAlSiO anorthite endmember CaAlSiO Solid solutions between K-feldspar and albite are called alkali feldspar. Solid solutions between albite and anorthite are called plagioclase, or, more properly, plagioclase feldspar. Only limited solid solution occurs between K-feldspar and anorthite, and in the two other solid solutions, immiscibility occurs at temperatures common in the crust of the Earth. Albite is considered both a plagioclase and alkali feldspar. The ratio of alkali feldspar to plagioclase feldspar, together with the proportion of quartz, is the basis for the QAPF classification of igneous rock. Calcium-rich plagioclase is the first feldspar to crystallize from cooling magma, then the plagioclase becomes increasingly sodium-rich as crystallization continues. This defines the continuous Bowen's reaction series. K-feldspar is the final feldspar to crystallize from the magma. Alkali feldspars Alkali feldspars are grouped into two types: those containing potassium in combination with sodium, aluminium, or silicon; and those where potassium is replaced by barium. The first of these include: orthoclase (monoclinic) sanidine (monoclinic) microcline (triclinic) anorthoclase (triclinic) Potassium and sodium feldspars are not perfectly miscible in the melt at low temperatures, therefore intermediate compositions of the alkali feldspars occur only in higher temperature environments. Sanidine is stable at the highest temperatures, and microcline at the lowest. Perthite is a typical texture in alkali feldspar, due to exsolution of contrasting alkali feldspar compositions during cooling of an intermediate composition. The perthitic textures in the alkali feldspars of many granites can be seen with the naked eye. Microperthitic textures in crystals are visible using a light microscope, whereas cryptoperthitic textures can be seen only with an electron microscope. Ammonium feldspar Buddingtonite is an ammonium feldspar with the chemical formula: NH4AlSi3O8. It is a mineral associated with hydrothermal alteration of the primary feldspar minerals. Barium feldspars Barium feldspars form as the result of the substitution of barium for potassium in the mineral structure. Barium feldspars are sometimes classified as a separate group of feldspars, and sometimes they are classified as a sub-group of alkali feldspars. The barium feldspars are monoclinic and include the following: celsian hyalophane Plagioclase feldspars The plagioclase feldspars are triclinic. The plagioclase series follows (with percent anorthite in parentheses): albite (0 to 10) oligoclase (10 to 30) andesine (30 to 50) labradorite (50 to 70) bytownite (70 to 90) anorthite (90 to 100) Intermediate compositions of exsolve to two feldspars of contrasting composition during cooling, but diffusion is much slower than in alkali feldspar, and the resulting two-feldspar intergrowths typically are too fine-grained to be visible with optical microscopes. The immiscibility gaps in the plagioclase solid solutions are complex compared to the gap in the alkali feldspars. The play of colours visible in some feldspar of labradorite composition is due to very fine-grained exsolution lamellae known as Bøggild intergrowth. The specific gravity in the plagioclase series increases from albite (2.62) to anorthite (2.72–2.75). Structure The structure of a feldspar crystal is based on aluminosilicate tetrahedra. Each tetrahedron consists of an aluminium or silicon ion surrounded by four oxygen ions. Each oxygen ion, in turn, is shared by a neighbouring tetrahedron to form a three-dimensional network. The structure can be visualized as long chains of aluminosilicate tetrahedra, sometimes described as crankshaft chains because their shape is kinked. Each crankshaft chain links to neighbouring crankshaft chains to form a three-dimensional network of fused four-member rings. The structure is open enough for cations (typically sodium, potassium, or calcium) to fit into the structure and provide charge balance. Weathering Chemical weathering of feldspars happens by hydrolysis and produces clay minerals, including illite, smectite, and kaolinite. Hydrolysis of feldspars begins with the feldspar dissolving in water, which happens best in acidic or basic solutions and less well in neutral ones. The speed at which feldspars are weathered is controlled by how quickly they are dissolved. Dissolved feldspar reacts with H+ or OH− ions and precipitates clays. The reaction also produces new ions in solution, with the variety of ions controlled by the type of feldspar reacting. The abundance of feldspars in the Earth's crust means that clays are very abundant weathering products. About 40% of minerals in sedimentary rocks are clays and clays are the dominant minerals in the most common sedimentary rocks, mudrocks. They are also an important component of soils. Feldspar that has been replaced by clay looks chalky compared to more crystalline and glassy unweathered feldspar grains. Feldspars, especially plagioclase feldspars, are not very stable at the Earth's surface due to their high formation temperature. This lack of stability is why feldspars are easily weathered to clays. Because of this tendency to weather easily, feldspars are usually not prevalent in sedimentary rocks. Sedimentary rocks that contain large amounts of feldspar indicate that the sediment did not undergo much chemical weathering before being buried. This means it was probably transported a short distance in cold and/or dry conditions that did not promote weathering, and that it was quickly buried by other sediment. Sandstones with large amounts of feldspar are called arkoses. Applications Feldspar is a common raw material used in glassmaking, ceramics, and to some extent as filler and an extender in paint, plastics, and rubber. In the US, about 66% of feldspar is consumed in glassmaking, including glass containers and glass fibre. Ceramics (including electrical insulators, sanitaryware, tableware and tile) and other uses, such as fillers, accounted for the remainder. Glass: Feldspar provides both K2O and Na2O for fluxing, and Al2O3 and CaO as stabilizers. As an important source of Al2O3 for glassmaking, feldspar is valued for its low iron and refractory mineral content, a low cost per unit of Al2O3, no volatiles and no waste. Ceramics: Feldspars are used in the ceramic industry as a flux to form a glassy phase in bodies during firing, and thus promote vitrification. They also are used as a source of alkalies and alumina in glazes. The composition of feldspar used in different ceramic formulations varies depending on various factors, including the properties of the individual grade, the other raw materials and the requirements of the finished products. However, typical additions include: tableware, 15% to 30% feldspar; high-tension electrical porcelains, 25% to 35%; sanitaryware, 25%; wall tile, 0% to 10%; and dental porcelain up to 80% feldspar. Earth sciences: In earth sciences and archaeology, feldspars are used for potassium-argon dating, argon-argon dating and luminescence dating. Minor use: Some household cleaners (such as Bar Keepers Friend and Bon Ami) use feldspar to give a mild abrasive action. Production The USGS estimated global production of feldspar in 2020 to be 26 million tonnes, with the top four producing countries being: China 2 million tonnes; India 5 million tonnes; Italy 4 million; Turkey 7.6 million tonnes. Commercial grades Typical mineralogical and chemical analyses of three commercial grades used in ceramics are: Extraterrestrial In October 2012, the Curiosity rover found high feldspar content in a Mars rock. Gallery	Physical sciences	Mineralogy	null
37977	https://en.wikipedia.org/wiki/Millet	Millet	Millets () are a highly varied group of small-seeded grasses, widely grown around the world as cereal crops or grains for fodder and human food. Most millets belong to the tribe Paniceae. Millets are important crops in the semiarid tropics of Asia and Africa, especially in India, Mali, Nigeria, and Niger, with 97% of production in developing countries. The crop is favoured for its productivity and short growing season under hot dry conditions. The millets are sometimes understood to include the widely cultivated sorghum; apart from that, pearl millet is the most commonly cultivated of the millets. Finger millet, proso millet, and foxtail millet are other important crop species. Millets may have been consumed by humans for about 7,000 years and potentially had "a pivotal role in the rise of multi-crop agriculture and settled farming societies". Etymology The name millet is derived via Old French millet, millot from Latin millium, 'millet', ultimately from Proto-Indo-European *mele-, 'to crush'. Description Characteristics Millets are small-grained, annual, warm-weather cereals belonging to the grass family. They are highly tolerant of drought and other extreme weather conditions and have a similar nutrient content to other major cereals. Taxonomic history In 1753, Carl Linnaeus described foxtail millet as Panicum italicum. In 1812, Palisot de Beauvois grouped several taxa into Setaria italica. The genus Pennisetum was divided by Otto Stapf in 1934 into the section penicillaria, with 32 species including all the cultivated ones, and four other sections. In 1977, J. Brunken and colleagues classed the wild P. violaceum as part of the cultivated species P. glaucum (pearl millet). Finger millet was described as Eleusine coracana by Joseph Gaertner in 1788. Evolution Phylogeny The millets are closely related to sorghum and maize within the PACMAD clade of grasses, and more distantly to the cereals of the BOP clade such as wheat and barley. Within the Panicoideae, sorghum is in the tribe Andropogoneae, while pearl millet, proso, foxtail, fonio, little millet, sawa, Japanese barnyard millet and kodo are in the tribe Paniceae. Within the Chloridoideae, finger millet is in the tribe Cynodonteae, while teff is in the tribe Eragrostideae. Taxonomy The different species of millets are not all closely related. All are members of the family Poaceae (the grasses), but they belong to different tribes and subfamilies. Commonly cultivated millets are: Eragrostideae tribe in the subfamily Chloridoideae: Eleusine coracana: Finger millet Eragrostis tef: Teff; often not considered to be a millet Paniceae tribe in the subfamily Panicoideae: Genus Panicum: Panicum miliaceum: Proso millet (common millet, broomcorn millet, hog millet, or white millet, also known as baragu in Kannada, panivaragu in Tamil) Panicum sumatrense: Little millet Panicum hirticaule: Sonoran millet, cultivated in the American Southwest Cenchrus americanus: Pearl millet Setaria italica: Foxtail millet, Italian millet, panic Genus Digitaria: of minor importance as crops Digitaria exilis: known as white fonio, fonio millet, and hungry rice or acha rice Digitaria iburua: Black fonio Digitaria compacta: Raishan, cultivated in the Khasi Hills of northeast India Digitaria sanguinalis: Polish millet Genus Echinochloa: collectively, the members of this genus are called barnyard grasses or barnyard millets Echinochloa esculenta: Japanese barnyard millet Echinochloa frumentacea: Indian barnyard millet Echinochloa stagnina: Burgu millet Echinochloa crus-galli: Common barnyard grass (or cockspur grass) Paspalum scrobiculatum: Kodo millet Genus Urochloa (formerly Brachiaria) Urochloa deflexa: Guinea millet Urochloa ramosa: Browntop millet, southern India Spodiopogon formosanus: Taiwan oil millet, endemic to Taiwan Andropogoneae tribe, also in the subfamily Panicoideae: Sorghum bicolor: Sorghum; usually considered a separate cereal, but sometimes known as great millet Coix lacryma-jobi: Job's tears, also known as adlay millet Domestication and spread Specialized archaeologists called palaeoethnobotanists, relying on data such as the relative abundance of charred grains found in archaeological sites, hypothesize that the cultivation of millets was of greater prevalence in prehistory than rice, especially in northern China and Korea. The cultivation of common millet as the earliest dry crop in East Asia has been attributed to its resistance to drought, and this has been suggested to have aided its spread. Asian varieties of millet made their way from China to the Black Sea region of Europe by 5000 BC. Millet was growing wild in Greece as early as 3000 BC, and bulk storage containers for millet have been found from the Late Bronze Age in Macedonia and northern Greece. Hesiod describes that "the beards grow round the millet, which men sow in summer." Millet is listed along with wheat in the third century BC by Theophrastus in his Enquiry into Plants. East Asia Proso millet (Panicum miliaceum) and foxtail millet (Setaria italica) were important crops beginning in the Early Neolithic of China. Some of the earliest evidence of millet cultivation in China was found at Cishan (north), where proso millet husk phytoliths and biomolecular components have been identified around 10,300–8,700 years ago in storage pits along with remains of pit-houses, pottery, and stone tools related to millet cultivation. Evidence at Cishan for foxtail millet dates back to around 8,700 years ago. Noodles made from these two varieties of millet were found under a 4,000-year-old earthenware bowl containing well-preserved noodles at the Lajia archaeological site in north China; this is the oldest evidence of millet noodles in China. Palaeoethnobotanists have found evidence of the cultivation of millet in the Korean Peninsula dating to the Middle Jeulmun pottery period (around 3500–2000 BC). Millet continued to be an important element in the intensive, multicropping agriculture of the Mumun pottery period (about 1500–300 BC) in Korea. Millets and their wild ancestors, such as barnyard grass and panic grass, were also cultivated in Japan during the Jōmon period sometime after 4000 BC. In the Zhengluo region of China, two millet species (foxtail millet and proso millet) were grown, enabling the people to survive the cooling of the global climate around 2200 BC. Chinese myths attribute the domestication of millet to Shennong, a legendary Emperor of China, and Hou Ji, whose name means Lord Millet. Indian subcontinent Little millet (Panicum sumatrense) is believed to have been domesticated around 5000 BC in Indian subcontinent and Kodo millet (Paspalum scrobiculatum) around 3700 BC, also in Indian subcontinent. Browntop millet (Urochloa ramosa) was likely domesticated in the Deccan near the beginning of the third millennium BCE and spread throughout India though was later superseded by other millets. Various millets have been mentioned in some of the Yajurveda texts, identifying foxtail millet (priyaṅgu), Barnyard millet (aṇu) and black finger millet (śyāmāka), indicating that millet cultivation was happening around 1200 BC in India. Upon request by the Indian Government in 2018, the Food and Agriculture Organisation of the United Nations declared 2023 as International Year of Millets. Africa Pearl millet (Pennisetum glaucum) was domesticated in the Sahel region of West Africa from Pennisetum violaceum. Early archaeological evidence in Africa includes finds at Birimi in northern Ghana (1740 cal BC) and Dhar Tichitt in Mauritania (1936–1683 cal BC) and the lower Tilemsi valley in Mali (2500 to 2000 cal BC). Studies of isozymes suggest domestication took place north east of the Senegal River in the far west of the Sahel and tentatively around 6000 BC. Pearl millet had arrived in the Indian subcontinent by 2000 BC to 1700 BC. Finger millet is native to the highlands of East Africa and was domesticated before the third millennium BC. Its cultivation had spread to South India by 1800 BC. Europe Broomcorn or proso millet (Panicum miliaceum) came to Europe from East Asia as early as the 17th century BC in Vinogradnyi Sad, Ukraine. At around 1500 BC it reached Italy and southeastern Europe; around 1400 BC it came to central Europe, and from 1200 BC, it arrived in northern Germany. Agriculture Cultivation Pearl millet is one of the two major crops in the semiarid, impoverished, less fertile agriculture regions of Africa and southeast Asia. Millets are not only adapted to poor, dry infertile soils, but they are also more reliable under these conditions than most other grain crops. Millets, however, do respond to high fertility and moisture. On a per-hectare basis, millet grain production can be 2 to 4 times higher with use of irrigation and soil supplements. Improved breeds of millet with enhanced disease resistance can significantly increase farm yield. There has been cooperation between poor countries to improve millet yields. For example, 'Okashana 1', a variety developed in India from a natural-growing millet variety in Burkina Faso, doubled yields. This breed was selected for trials in Zimbabwe. From there it was taken to Namibia, where it was released in 1990 and enthusiastically adopted by farmers. 'Okashana 1' became the most popular variety in Namibia, the only non-Sahelian country where pearl millet—locally known as mahangu—is the dominant food staple for consumers. 'Okashana 1' was then introduced to Chad. The breed has significantly enhanced yields in Mauritania and Benin. Pests and diseases Millets are subject to damage by many insect pests, including corn borers, stem borers, the caterpillars of numerous moths in the families Erebidae and Noctuidae, the millet midge, many species of flies in the Muscidae, Hemipteran bugs of many families including aphids, and species of thrips, beetles, and grasshoppers. Among the many diseases of millets are serious fungal infections such as anthracnose, blast, charcoal rot, downy mildew, ergot, grain mould, rust, and sheath rot. Bacterial diseases are generally less serious; they include bacterial leaf spot, leaf stripe and leaf streak. Viral diseases are again generally less serious, except for a few diseases such as maize stripe virus, maize mosaic virus, sorghum red stripe virus, and maize streak virus. Production In 2022, global production of millet was 30.9 million tonnes. India is the top millet producer worldwide, with 11.8 million tonnes grown annually – some 38% of the world total and nearly triple its nearest rival. Eight of the remaining nine nations in the top 10 producers are in Africa, ranging from Niger (at 3.7 million tonnes) to Chad (0.7 million tonnes); the sole exception is China, number three in global production, at 2.7 million tonnes. Research Research on millets is carried out by the International Crops Research Institute for the Semi-Arid Tropics (ICRISAT) and ICAR-Indian Institute of Millets Research in Telangana, India, and by the United States Department of Agriculture's Agricultural Research Service at Tifton, Georgia, United States. Uses As food In Ukraine, millet was historically a common ingredient in the diet of the Zaporozhian Cossacks, in the form of a porridge called "kulish". This dish, primarily made with millet, served with stewed vegetables and meat, cooked in a cauldron, remains a part of modern Ukrainian cuisine. In Germany, it is eaten sweet, for example with milk and berries for breakfast. Millet is the main ingredient in , a Vietnamese sweet snack. It contains a layer of smashed millet and mungbean topped with sliced dried coconut meat wrapped in a crunchy rice cake. In parts of Africa millet is mixed with milk to make a drink, Brukina. Finger millet is made into ragi rotti flatbread and ragi mudde dough lumps in Karnataka. Dough lumps are eaten as fura in the Sahel region of West Africa. Alcoholic beverages In the Himalayas, including in Nepal, Sikkim, and Darjeeling, millet is fermented into Tongba, an alcoholic drink. In India, alcoholic beverages including rakshi are produced from millets. As forage Millet is sometimes used as a forage crop. Compared to forage sorghum, animals including lambs gain weight faster on millet, and it has better hay or silage potential, although it produces less dry matter. Millet does not contain toxic prussic acid, sometimes found in sorghum. The rapid growth of millet as a grazing crop allows flexibility in its use. Farmers can wait until sufficient late spring / summer moisture is present and then make use of it. It is ideally suited to irrigation where livestock finishing is required. Human consumption Per capita consumption of millets as food varies in different parts of the world, with consumption being the highest in Western Africa. In the Sahel region, millet is estimated to account for about 35 percent of total cereal food consumption in Burkina Faso, Chad and the Gambia. In Mali and Senegal, millets constitute roughly 40 percent of total cereal food consumption per capita, while in Niger and arid Namibia it is over 65 percent (see mahangu). Other countries in Africa where millets are a significant food source include Ethiopia, Nigeria and Uganda. Millet is also an important food item for the population living in the drier parts of many other countries, especially in eastern and central Africa, and in the northern coastal countries of western Africa. In developing countries outside Africa, millet has local significance as a food in parts of some countries, such as China, India, Burma and North Korea. People affected by gluten-related disorders, such as coeliac disease, non-celiac gluten sensitivity and wheat allergy sufferers, who need a gluten-free diet, can replace gluten-containing cereals in their diets with millet. There remains a risk of contamination with gluten-containing cereals. Nutrition The table shows the nutrient content of the grains of different species of millet, raw, compared to other staples.	Biology and health sciences	Poales	null
37993	https://en.wikipedia.org/wiki/Toilet%20paper	Toilet paper	Toilet paper (sometimes called toilet/bath/bathroom tissue, or toilet roll) is a tissue paper product primarily used to clean the anus and surrounding region of feces (after defecation), and to clean the external genitalia and perineal area of urine (after urination). It is commonly supplied as a long strip of perforated paper wrapped around a cylindrical paperboard core, for storage in a dispenser within arm's reach of a toilet. The bundle, or roll of toilet paper, is specifically known as a toilet roll, loo roll, or bog roll (in Britain). There are other uses for toilet paper, as it is a readily available household product. It can be used for blowing the nose or wiping the eyes (or other uses of facial tissue). It can be used to wipe off sweat or absorb it. Some people may use the paper to absorb the bloody discharge that comes out of the vagina during menstruation. Toilet paper can be used in cleaning (like a less abrasive paper towel). As a teenage prank, "toilet papering" is a form of temporary vandalism. Most modern toilet paper in the developed world is designed to decompose in septic tanks, whereas some other bathroom and facial tissues are not. Wet toilet paper rapidly decomposes in the environment. Toilet paper comes in various numbers of plies (layers of thickness), from one- to six-ply, with more back-to-back plies providing greater strength and absorbency. Most modern domestic toilet paper is white, and patterned or textured. Some people have a preference for whether the orientation of the roll on a dispenser should be over or under. The use of paper for hygiene has been recorded in China in the 6th century AD, with specifically manufactured toilet paper being mass-produced in the 14th century. Modern commercial toilet paper originated in the 19th century, with a patent for roll-based dispensers being made in 1883. History Although paper had been known as a wrapping and padding material in China since the 2nd century BC, a reference to the use of toilet paper dates back as early as when the scholar-official Yan Zhitui (531–591) wrote: During the later Tang dynasty (618–907 AD), an Arab traveller to China in the year 851 AD remarked: During the early 14th century, it was recorded that in what is now Zhejiang alone, ten million packages of 1,000 to 10,000 sheets of toilet paper were manufactured annually. During the Ming dynasty (1368–1644 AD), it was recorded in 1393 that an annual supply of 720,000 sheets of toilet paper (approximately ) were produced for the general use of the imperial court at the capital of Nanjing. From the records of the Imperial Bureau of Supplies of that same year, it was also recorded that for the Hongwu Emperor's imperial family alone, there were 15,000 sheets of special soft-fabric toilet paper made, and each sheet of toilet paper was perfumed. Elsewhere, wealthy people wiped themselves with wool, lace or hemp, while less wealthy people used their hand when defecating into rivers, or cleaned themselves with various materials such as rags, wood shavings, leaves, grass, hay, stones, sand, moss, water, snow, ferns, plant husks, fruit skins, seashells, or corncobs, depending upon the country and weather conditions or social customs. In Ancient Rome, a sponge on a stick was commonly used, and, after use, placed back in a pail of vinegar. Several talmudic sources indicating ancient Jewish practice refer to the use of small pebbles, often carried in a special bag, and also to the use of dry grass and of the smooth edges of broken pottery jugs (e.g., Shabbat 81a, 82a, Yevamot 59b). These are all cited in the classic Biblical and Talmudic Medicine by the German physician Julius Preuss (Eng. trans. Sanhedrin Press, 1978). The 16th-century French satirical writer François Rabelais, in Chapter XIII of Book 1 of his novel sequence Gargantua and Pantagruel, has his character Gargantua investigate a great number of ways of cleansing oneself after defecating. Gargantua dismisses the use of paper as ineffective, rhyming that: "Who his foul tail with paper wipes, Shall at his ballocks leave some chips." (Sir Thomas Urquhart's 1653 English translation). He concludes that "the neck of a goose, that is well downed" provides an optimum cleansing medium. The rise of publishing by the eighteenth century led to the use of newspapers and cheap editions of popular books for cleansing. Lord Chesterfield, in a letter to his son in 1747, told of a man who purchased In many parts of the world, especially where toilet paper or the necessary plumbing for disposal may be unavailable or unaffordable, toilet paper is not used. Also, in many parts of the world people consider using water a much cleaner and more sanitary practice than using paper. Cleansing is then performed with other methods or materials, such as water, for example using a bidet, a lota, rags, sand, leaves (including seaweed), corn cobs, animal furs, sticks or hands; afterwards, hands are washed with water and possibly soap. On 18 July 2024 the sale of ruble-note artwork on toilet paper was banned by a Moscow judge. As a commodity Joseph Gayetty is widely credited with being the inventor of modern commercially available toilet paper in the United States. Gayetty's paper, first introduced in 1857, was available as late as the 1920s. Gayetty's Medicated Paper was sold in packages of flat sheets, watermarked with the inventor's name. Original advertisements for the product used the tagline "The greatest necessity of the age! Gayetty's medicated paper for the water-closet". Seth Wheeler of Albany, New York, obtained the earliest United States patents for toilet paper and dispensers, the types of which eventually were in common use in that country, in 1883. Toilet paper dispensed from rolls was popularized when the Scott Paper Company began marketing it in 1890. The manufacturing of this product had a long period of refinement, considering that as late as the 1930s, a selling point of the Northern Tissue company was that their toilet paper was "splinter free". The widespread adoption of the flush toilet increased the use of toilet paper, as heavier paper was more prone to clogging the trap that prevents sewer gases from escaping through the toilet. Softer, two ply toilet roll was introduced in Britain in 1942, by St Andrew Mills in Walthamstow; this became the famous Andrex. Moist toilet paper, called wet wipes, was first introduced in the United Kingdom by Andrex in the 1990s. It has been promoted as being a better method of cleaning than dry toilet paper after defecation, and may be useful for women during menstruation. It was promoted as a flushable product but it has been implicated in the creation of fatbergs; by 2016 some municipalities had begun education campaigns advising people not to flush used wet wipes. More than seven billion rolls of toilet paper are sold yearly in the United States where an average of 23.6 rolls per capita per year is used. In 1973, Johnny Carson joked in his Tonight Show monologue about comments made by Wisconsin congressman Harold V. Froehlich about the possibility of a toilet paper shortage. Subsequently, consumers purchased abnormal amounts, causing an actual shortage in the United States for several months. Toilet paper has been one of the commodities subject to shortages in Venezuela starting in the 2010s; the government seized one toilet paper factory in an effort to resolve the problem. During the COVID-19 pandemic, toilet paper shortages were reported in March 2020 in multiple countries due to hoarding and panic buying. At first, few believed the pandemic would be serious. Later, people realized they might need to stock up on certain items in case of a shelter-in-place order, or in case they did not know how long such an order would last; suppliers could not assure that they could keep up with demand. However, manufacturers continued to produce even more than they had before. Demand was higher for the types of toilet paper used at home. In some countries the bidet was already seen as a solution, and a survey before the pandemic had indicated an increasing number of Americans would be interested. Amid the panic buying during the pandemic, the Australian toilet paper brand Quilton donated a million of toilet paper rolls to vulnerable Australians who were struggling due to the shortages of toilet paper. In 2022, British toilet paper packaging started displaying bowel cancer symptoms to raise awareness, following campaigning from blogger and journalist Deborah James, who later died from the disease in June 2022. At the time, half of all Britons could not name any of the main symptoms of bowel cancer. Andrex were the first brand to take the lead on the matter, then various supermarkets followed suit. Description Toilet paper is available in several types of paper, a variety of patterns, decorations, and textures, and it may be moistened or perfumed, although fragrances sometimes cause problems for users who are allergic to perfumes. The average measures of a modern roll of toilet paper is c. 10 cm (3 in.) wide, and 12 cm (4 in.) in diameter, and weighs between and . An alternative method of packing the sheets uses interleaved sheets in boxes, or in bulk for use in dispensers. "Hard" single-ply paper has been used as well as soft multi-ply. Sheet size The format of individual sheets of toilet paper, which is given by a perforation line, varies nationally. In Germany, Holland, France, Poland, Switzerland, for example, about postcard size is standard (about 100 × 140 mm), so about DIN format (DIN A6 105 × 148 mm). In England, the usual format is already somewhat wider, about 115 × 135 mm. The most extreme landscape format with 115 × 102 mm exists in Thailand. The most extreme portrait format (not counting toilet paper rolls without any perforation) is 100 × 366 mm; a promotional toilet paper from Schmidt Spiele in Germany. Manufactured toilet paper sheet in the United States was sized × . Since 1999 the size of a sheet has been shrinking; Kimberly-Clark reduced the length of a sheet to . Scott, in 2006, reduced the length of their product to . The width of sheets was later reduced giving a general sheet size of long and wide. Larger sizes remain available. Sheet ply The ply of a toilet paper refers to the number of layers per sheet. Rolls are typically available in single-ply, 2-ply, 3-ply, and 4-ply. Roll length Phrases like "single roll", "double roll", "triple roll", "jumbo roll", and "mega roll" commonly used in retail advertising refer to the number of sheets per roll (though the actual number of sheets is also usually disclosed on packaging). A longer roll needs to be replaced less often, but the very largest sizes do not fit all toilet paper dispensers, especially in older homes. Materials Toilet paper is usually manufactured from pulpwood trees, but is also sometimes made from sugar cane byproducts or bamboo. Toilet paper products vary greatly in the distinguishing technical factors, such as size, weight, roughness, softness, chemical residues, "finger-breakthrough" resistance, water-absorption, etc. The larger companies have very detailed, scientific market surveys to determine which marketing sectors require or demand which of the many technical qualities. Modern toilet paper may have a light coating of aloe or lotion or wax worked into the paper to reduce roughness. Quality is usually determined by the number of plies (stacked sheets), coarseness, and durability. Low grade institutional toilet paper is typically of the lowest grade of paper, has only one or two plies, is very coarse and sometimes contains small amounts of embedded unbleached/unpulped paper; it was typically called "hard" toilet paper. A brand disinfected with carbolic acid was manufactured in Sheffield, United Kingdom under the Izal brand name by Newton Chambers until 1981. Mid-grade two ply is somewhat textured to provide some softness and is somewhat stronger. Premium toilet paper may have lotion and wax and has two to four plies of very finely pulped paper. If it is marketed as "luxury", it may be quilted or rippled (embossed), perfumed, colored or patterned, medicated (with anti-bacterial chemicals), or treated with aloe or other perfumes. To advance decomposition of the paper in septic tanks or drainage, the paper used has shorter fibres than facial tissue or writing paper. The manufacturer tries to reach an optimal balance between rapid decomposition (which requires shorter fibres) and sturdiness (which requires longer fibres). Compaction of toilet paper in drain lines, such as in a clog, prevents fibre dispersion and largely halts the breakdown process. A German quip says that the toilet paper of Nazi Germany was so rough and scratchy that it was almost unusable, so many people used old issues of the Völkischer Beobachter instead, because the paper was softer. Color and design Colored toilet paper in colors such as pink, lavender, light blue, light green, purple, green, and light yellow (so that one could choose a color of toilet paper that matched or complemented the color of one's bathroom) was commonly sold in the United States from the 1960s. Up until 2004, Scott was one of the last remaining U.S. manufacturers to still produce toilet paper in beige, blue, and pink. However, the company has since cut production of colored paper altogether. Colored toilet paper remains commonly available in some European countries. Here in solid color toilet paper base, apart from the natural tones between white and gray or beige, pastel shades prevail: pink, apricot, light yellow and light blue. In rare cases, pale purple or pale green can be found. However, rich colors are rarely used, such as black, wine red, neon green, royal blue. Flat printed toilet paper is uncommon. If there is printing, it is often one color. Common print colors are pink and pinkish red, also blue, more rarely purple, orange, brown or green. Design Today, in the United States, plain unpatterned colored toilet paper has been mostly replaced by patterned toilet paper, normally white, with embossed decorative patterns or designs in various colors and different sizes depending on the brand. The patterns are in most cases "scatter patterns", that is, a motif is distributed ("scattered") several times (irregularly) over the surface. Stripes and dot patterns are rare. Occasionally, toilet papers have an embossed crocodile, wave, circle or check pattern. Some are additionally printed. Ornaments usually stand on their own as self-contained units. They never go uninterrupted (for example, as a border) from the first to the last sheet. Motifs Predominant is everything that is associated "softness" and "fluffiness". There are decorations with bears, cats, rabbits, down feathers, clouds. Another motifs are things associated with "lightness": Clouds, downy feathers, leaves of all kinds, butterflies, flying birds. Another association is anything associated with pleasant fragrance: especially flowers of all kinds. Rare are motifs intended to appear noble, such as the Bourbon lily. Less rare are allusions to water, such as fish, shells and other aquatic creatures. Toilet papers are also provided with texts (jokes, poems), joke motifs (banknotes) or advertising imprints. Texture Toilet paper is offered in different qualities. The cheapest toilet papers have a texture close to crêpe paper. They are often made of recycled material. Expensive toilet papers are made from particularly absorbent, delicate tissue paper. Toilet paper usually has a smooth surface. With several intentions, it is occasionally embossed. On the one hand, the embossing can serve to stabilize the paper. Furthermore, wiping can become more effective. Thirdly, there are design reasons. In Switzerland, in particular, there are often toilet paper with burls. In Germany, the number of plies is considered a quality feature. In the USA, Great Britain and Japan, the quality feature is that the toilet paper is as delicate and fine as possible. Additives Some toilet papers are perfumed. Popular scents are chamomile, peach or rose. Other toilet papers are impregnated with antibacterial additives. Installation Dispensers A toilet roll holder, also known as a toilet paper dispenser, is an item that holds a roll of toilet paper. There are at least seven types of holders: A horizontal piece of wire mounted on a hinge, hanging from a door or wall. A horizontal axle recessed in the wall. A vertical axle recessed in the wall A horizontal axle mounted on a freestanding frame. A freestanding vertical pole on a base. A wall mounted dispensing unit, usually containing more than one roll. This is used in the commercial/away-from-home marketplace. A wall mounted dispensing unit with tissue interleaved in a "S"-type fold so the user can extract the tissue one sheet at a time. Some commercial or institutional toilet paper is wrapped around a cylinder to many times the thickness of a standard toilet paper roll. Orientation There are two choices of orientation when using a holder with a horizontal axle parallel to the wall: the toilet paper may hang over or under the roll. The choice is largely a matter of personal preference, dictated by habit. In surveys of American consumers and of bath and kitchen specialists, 60–70% of respondents prefer over. Most Americans think it should go over the top, like a waterfall. Decoration Toilegami refers to toilet paper origami. Like table napkins, some fancy Japanese hotels fold the first squares of toilet paper on its dispenser to be presented in a fashionable way. Recreational use In the United States, toilet paper has been the primary tool in a prank known as "TP-ing" (pronounced "teepeeing"). TP-ing, or "toilet papering", is often favored by adolescents and is the act of throwing rolls of toilet paper over cars, trees, houses and gardens, causing the toilet paper to unfurl and cover the property, creating an inconvenient mess. Children and cats may unroll an entire roll of toilet paper by spinning it until it completely unravels on the floor, or as a game by children wadding up one end, putting it in the toilet bowl without tearing it and then using the flushing of the toilet to pull new paper into the toilet, with the objective of flushing the entire roll down the toilet section at a time without the toilet paper breaking. Special toilet paper insert holders with an oblong shape were invented to prevent continuous unrolling without tearing to discourage this practice. Toilet paper pranks include musical toilet paper holders and inserts that are activated by the unrolling of the toilet paper and will loudly play an embarrassing song calling attention to the person defecating. Other gags include custom toilet paper printed with jokes, stories or politician's images. Mechanics Alexander Balankin and coauthors have studied the behavior of toilet paper under tensile stress and during wetting and burning. Toilet paper has been used in physics education to demonstrate the concepts of torque, moment of inertia, and angular momentum; and the conservation of momentum and energy. Environmental considerations One tree produces about 800 rolls () of toilet paper and about 83 million rolls are produced per day. Global toilet paper production consumes 27,000 trees daily. More than seven billion rolls of toilet paper are sold yearly in the United States alone. Americans use an average of 141 rolls per capita a year which is equivalent to of tissue paper per year. This figure is about 50% more than the average of other Western countries or Japan. The higher use in the United States may be explained by the fact that other countries people use bidets or spray hoses to clean themselves. Millions of trees are harvested in North and South America leaving ecological footprint concerns. , between 22% and 48% of the toilet paper used in the United States comes from tree farms in the U.S. and South America, with the rest mostly coming from old, second growth forests, and, some from virgin forests. Alternatives to virgin wood pulp Toilet paper made from recycled paper avoids the direct environmental impact of cutting down trees, and is commercially available. Recycled newspaper can contain BPA, an endocrine disruptor. Toilet paper produced from bamboo is commercially available, and is in some ways more environmentally friendly than virgin pulpwood, because bamboo grows faster, taking less land and less water. For North American consumers, the Natural Resources Defense Council recommends recycled tree pulp over bamboo toilet paper, because tree forests promote more biodiversity and bamboo products must be shipped from Asia. Toilet paper produced from bagasse, a byproduct of sugarcane, is commercially available, and avoids cutting down any plants because sugarcane is already grown for sugar production. The most environmentally friendly alternatives are to rely solely on soap and water for anal hygiene.	Biology and health sciences	Hygiene products	Health
38001	https://en.wikipedia.org/wiki/Phenylalanine	Phenylalanine	Phenylalanine (symbol Phe or F) is an essential α-amino acid with the formula . It can be viewed as a benzyl group substituted for the methyl group of alanine, or a phenyl group in place of a terminal hydrogen of alanine. This essential amino acid is classified as neutral, and nonpolar because of the inert and hydrophobic nature of the benzyl side chain. The L-isomer is used to biochemically form proteins coded for by DNA. Phenylalanine is a precursor for tyrosine, the monoamine neurotransmitters dopamine, norepinephrine (noradrenaline), and epinephrine (adrenaline), and the biological pigment melanin. It is encoded by the messenger RNA codons UUU and UUC. Phenylalanine is found naturally in the milk of mammals. It is used in the manufacture of food and drink products and sold as a nutritional supplement as it is a direct precursor to the neuromodulator phenethylamine. As an essential amino acid, phenylalanine is not synthesized de novo in humans and other animals, who must ingest phenylalanine or phenylalanine-containing proteins. The one-letter symbol F was assigned to phenylalanine for its phonetic similarity. History The first description of phenylalanine was made in 1879, when Schulze and Barbieri identified a compound with the empirical formula, C9H11NO2, in yellow lupine (Lupinus luteus) seedlings. In 1882, Erlenmeyer and Lipp first synthesized phenylalanine from phenylacetaldehyde, hydrogen cyanide, and ammonia. The genetic codon for phenylalanine was first discovered by J. Heinrich Matthaei and Marshall W. Nirenberg in 1961. They showed that by using mRNA to insert multiple uracil repeats into the genome of the bacterium E. coli, they could cause the bacterium to produce a polypeptide consisting solely of repeated phenylalanine amino acids. This discovery helped to establish the nature of the coding relationship that links information stored in genomic nucleic acid with protein expression in the living cell. Dietary sources Good sources of phenylalanine are eggs, chicken, liver, beef, milk, and soybeans. Another common source of phenylalanine is anything sweetened with the artificial sweetener aspartame, such as diet drinks, diet foods and medication; the metabolism of aspartame produces phenylalanine as one of the compound's metabolites. Dietary recommendations The Food and Nutrition Board (FNB) of the U.S. Institute of Medicine set Recommended Dietary Allowances (RDAs) for essential amino acids in 2002. For phenylalanine plus tyrosine, for adults 19 years and older, 33 mg/kg body weight/day. In 2005 the DRI is set to 27 mg/kg per day (with no tyrosine), the FAO/WHO/UNU recommendation of 2007 is 25 mg/kg per day (with no tyrosine). Other biological roles L-Phenylalanine is biologically converted into L-tyrosine, another one of the DNA-encoded amino acids. L-tyrosine in turn is converted into L-DOPA, which is further converted into dopamine, norepinephrine (noradrenaline), and epinephrine (adrenaline). The latter three are known as the catecholamines. Phenylalanine uses the same active transport channel as tryptophan to cross the blood–brain barrier. In excessive quantities, supplementation can interfere with the production of serotonin and other aromatic amino acids as well as nitric oxide due to the overuse (eventually, limited availability) of the associated cofactors, iron or tetrahydrobiopterin. The corresponding enzymes for those compounds are the aromatic amino acid hydroxylase family and nitric oxide synthase. In plants Phenylalanine is the starting compound used in the synthesis of flavonoids. Lignan is derived from phenylalanine and from tyrosine. Phenylalanine is converted to cinnamic acid by the enzyme phenylalanine ammonia-lyase. Biosynthesis Phenylalanine is biosynthesized via the shikimate pathway. Phenylketonuria The genetic disorder phenylketonuria (PKU) is the inability to metabolize phenylalanine because of a lack of the enzyme phenylalanine hydroxylase. Individuals with this disorder are known as "phenylketonurics" and must regulate their intake of phenylalanine. Phenylketonurics often use blood tests to monitor the amount of phenylalanine in their blood. Lab results may report phenylalanine levels using either mg/dL and μmol/L. One mg/dL of phenylalanine is approximately equivalent to 60 μmol/L. A (rare) "variant form" of phenylketonuria called hyperphenylalaninemia is caused by the inability to synthesize a cofactor called tetrahydrobiopterin, which can be supplemented. Pregnant women with hyperphenylalaninemia may show similar symptoms of the disorder (high levels of phenylalanine in blood), but these indicators will usually disappear at the end of gestation. Pregnant women with PKU must control their blood phenylalanine levels even if the fetus is heterozygous for the defective gene because the fetus could be adversely affected due to hepatic immaturity. A non-food source of phenylalanine is the artificial sweetener aspartame. This compound is metabolized by the body into several chemical byproducts including phenylalanine. The breakdown problems phenylketonurics have with the buildup of phenylalanine in the body also occurs with the ingestion of aspartame, although to a lesser degree. Accordingly, all products in Australia, the U.S. and Canada that contain aspartame must be labeled: "Phenylketonurics: Contains phenylalanine." In the UK, foods containing aspartame must carry ingredient panels that refer to the presence of "aspartame or E951" and they must be labeled with a warning "Contains a source of phenylalanine." In Brazil, the label "Contém Fenilalanina" (Portuguese for "Contains Phenylalanine") is also mandatory in products which contain it. These warnings are placed to help individuals avoid such foods. D-, L- and DL-phenylalanine The stereoisomer D-phenylalanine (DPA) can be produced by conventional organic synthesis, either as a single enantiomer or as a component of the racemic mixture. It does not participate in protein biosynthesis although it is found in proteins in small amounts - particularly aged proteins and food proteins that have been processed. The biological functions of D-amino acids remain unclear, although D-phenylalanine has pharmacological activity at niacin receptor 2. DL-Phenylalanine (DLPA) is marketed as a nutritional supplement for its purported analgesic and antidepressant activities, which have been supported by clinical trials. DL-Phenylalanine is a mixture of D-phenylalanine and L-phenylalanine. The reputed analgesic activity of DL-phenylalanine may be explained by the possible blockage by D-phenylalanine of enkephalin degradation by the enzyme carboxypeptidase A. Enkephalins act as agonists of the mu and delta opioid receptors, and agonists of these receptors are known to produce antidepressant effects. The mechanism of DL-phenylalanine's supposed antidepressant activity may also be accounted for in part by the precursor role of L-phenylalanine in the synthesis of the neurotransmitters norepinephrine and dopamine, though clinical trials have not found an antidepressant effect from L-phenylalanine alone. Elevated brain levels of norepinephrine and dopamine are thought to have an antidepressant effect. D-Phenylalanine is absorbed from the small intestine and transported to the liver via the portal circulation. A small amount of D-phenylalanine appears to be converted to L-phenylalanine. D-Phenylalanine is distributed to the various tissues of the body via the systemic circulation. It appears to cross the blood–brain barrier less efficiently than L-phenylalanine, and so a small amount of an ingested dose of D-phenylalanine is excreted in the urine without penetrating the central nervous system. L-Phenylalanine is an antagonist at α2δ Ca2+ calcium channels with a Ki of 980 nM. In the brain, L-phenylalanine is a competitive antagonist at the glycine binding site of NMDA receptor and at the glutamate binding site of AMPA receptor. At the glycine binding site of NMDA receptor L-phenylalanine has an apparent equilibrium dissociation constant (KB) of 573 μM estimated by Schild regression which is considerably lower than brain L-phenylalanine concentration observed in untreated human phenylketonuria. L-Phenylalanine also inhibits neurotransmitter release at glutamatergic synapses in hippocampus and cortex with IC50 of 980 μM, a brain concentration seen in classical phenylketonuria, whereas D-phenylalanine has a significantly smaller effect. Commercial synthesis L-Phenylalanine is produced for medical, feed, and nutritional applications, such as aspartame, in large quantities by utilizing the bacterium Escherichia coli, which naturally produces aromatic amino acids like phenylalanine. The quantity of L-phenylalanine produced commercially has been increased by genetically engineering E. coli, such as by altering the regulatory promoters or amplifying the number of genes controlling enzymes responsible for the synthesis of the amino acid. Derivatives Boronophenylalanine (BPA) is a dihydroxyboryl derivative of phenylalanine, used in neutron capture therapy. 4-Azido-L-phenylalanine is a protein-incorporated unnatural amino acid used as a tool for bioconjugation in the field of chemical biology.	Biology and health sciences	Amino acids	Biology
38011	https://en.wikipedia.org/wiki/Squid	Squid	A squid (: squid) is a mollusc with an elongated soft body, large eyes, eight arms, and two tentacles in the orders Myopsida, Oegopsida, and Bathyteuthida (though many other molluscs within the broader Neocoleoidea are also called squid despite not strictly fitting these criteria). Like all other cephalopods, squid have a distinct head, bilateral symmetry, and a mantle. They are mainly soft-bodied, like octopuses, but have a small internal skeleton in the form of a rod-like gladius or pen, made of chitin. Squid diverged from other cephalopods during the Jurassic and occupy a similar role to teleost fish as open water predators of similar size and behaviour. They play an important role in the open water food web. The two long tentacles are used to grab prey and the eight arms to hold and control it. The beak then cuts the food into suitable size chunks for swallowing. Squid are rapid swimmers, moving by jet propulsion, and largely locate their prey by sight. They are among the most intelligent of invertebrates, with groups of Humboldt squid having been observed hunting cooperatively. They are preyed on by sharks, other fish, sea birds, seals and cetaceans, particularly sperm whales. Squid can change colour for camouflage and signalling. Some species are bioluminescent, using their light for counter-illumination camouflage, while many species can eject a cloud of ink to distract predators. Squid are used for human consumption with commercial fisheries in Japan, the Mediterranean, the southwestern Atlantic, the eastern Pacific and elsewhere. They are used in cuisines around the world, often known as "calamari". Squid have featured in literature since classical times, especially in tales of giant squid and sea monsters. Taxonomy and phylogeny Squid are members of the class Cephalopoda, subclass Coleoidea. The squid orders Myopsida and Oegopsida are in the superorder Decapodiformes (from the Greek for "ten-legged"). Two other orders of decapodiform cephalopods are also called squid, although they are taxonomically distinct from squids and differ recognizably in their gross anatomical features. They are the bobtail squid of order Sepiolida and the ram's horn squid of the monotypic order Spirulida. The vampire squid (Vampyroteuthis infernalis), however, is more closely related to the octopus than to any squid. The cladogram, not fully resolved, is based on Sanchez et al., 2018. Their molecular phylogeny used mitochondrial and nuclear DNA marker sequences; they comment that a robust phylogeny "has proven very challenging to obtain". If it is accepted that Sepiidae cuttlefish are a kind of squid, then the squids, excluding the vampire squid, form a clade as illustrated. Orders are shown in boldface; all the families not included in those orders are in the paraphyletic order "Oegopsida", except Sepiadariidae and Sepiidae that are in the polyphyletic order "Sepiida", Evolution Crown coleoids (the common ancestor of octopuses and squid) diverged in the late Paleozoic (Mississippian), according to fossils of Syllipsimopodi, an early relative of vampire squids and octopuses. True squid diverged during the Jurassic, but many squid families appeared in or after the Cretaceous. Both the coleoids and the teleost fish were involved in much adaptive radiation at this time, and the two modern groups resemble each other in size, ecology, habitat, morphology and behaviour, however some fish moved into fresh water while the coleoids remained in marine environments. The ancestral coleoid was probably nautiloid-like with a strait septate shell that became immersed in the mantle and was used for buoyancy control. Four lines diverged from this, Spirulida (with one living member), the cuttlefishes, the squids and the octopuses. Squid have differentiated from the ancestral mollusc such that the body plan has been condensed antero-posteriorly and extended dorso-ventrally. What may have been the foot of the ancestor is modified into a complex set of appendages around the mouth. The sense organs are highly developed and include advanced eyes similar to those of vertebrates. The ancestral shell has been lost, with only an internal gladius, or pen, remaining. The pen, made of a chitin-like material, is a feather-shaped internal structure that supports the squid's mantle and serves as a site for muscle attachment. The cuttlebone or sepion of the Sepiidae is calcareous and appears to have evolved afresh in the Tertiary. Description Squid are soft-bodied molluscs whose forms evolved to adopt an active predatory lifestyle. The head and foot of the squid are at one end of a long body, and this end is functionally anterior, leading the animal as it moves through the water. A set of eight arms and two distinctive tentacles surround the mouth; each appendage takes the form of a muscular hydrostat and is flexible and prehensile, usually bearing disc-like suckers. The suckers may lie directly on the arm or be stalked. Their rims are stiffened with chitin and may contain minute toothlike denticles. These features, as well as strong musculature, and a small ganglion beneath each sucker to allow individual control, provide a very powerful adhesion to grip prey. Hooks are present on the arms and tentacles in some species, but their function is unclear. The two tentacles are much longer than the arms and are retractile. Suckers are limited to the spatulate tip of the tentacle, known as the manus. In the mature male, the outer half of one of the left arms is hectocotylised – and ends in a copulatory pad rather than suckers. This is used for depositing a spermatophore inside the mantle cavity of a female. A ventral part of the foot has been converted into a funnel through which water exits the mantle cavity. The main body mass is enclosed in the mantle, which has a swimming fin along each side. These fins are not the main source of locomotion in most species. The mantle wall is heavily muscled and internal. The visceral mass, which is covered by a thin, membranous epidermis, forms a cone-shaped posterior region known as the "visceral hump". The mollusc shell is reduced to an internal, longitudinal chitinous "pen" in the functionally dorsal part of the animal; the pen acts to stiffen the squid and provides attachments for muscles. On the functionally ventral part of the body is an opening to the mantle cavity, which contains the gills (ctenidia) and openings from the excretory, digestive and reproductive systems. An inhalant siphon behind the funnel draws water into the mantle cavity via a valve. The squid uses the funnel for locomotion via precise jet propulsion. In this form of locomotion, water is sucked into the mantle cavity and expelled out of the funnel in a fast, strong jet. The direction of travel is varied by the orientation of the funnel. Squid are strong swimmers and certain species can "fly" for short distances out of the water. Camouflage Squid make use of different kinds of camouflage, namely active camouflage for background matching (in shallow water) and counter-illumination. This helps to protect them from their predators and allows them to approach their prey. The skin is covered in controllable chromatophores of different colours, enabling the squid to match its coloration to its surroundings. The play of colours may in addition distract prey from the squid's approaching tentacles. The skin also contains light reflectors called iridophores and leucophores that, when activated, in milliseconds create changeable skin patterns of polarized light. Such skin camouflage may serve various functions, such as communication with nearby squid, prey detection, navigation, and orientation during hunting or seeking shelter. Neural control of the iridophores enabling rapid changes in skin iridescence appears to be regulated by a cholinergic process affecting reflectin proteins. Some mesopelagic squid such as the firefly squid (Watasenia scintillans) and the midwater squid (Abralia veranyi) use counter-illumination camouflage, generating light to match the downwelling light from the ocean surface. This creates the effect of countershading, making the underside lighter than the upperside. Counter-illumination is also used by the Hawaiian bobtail squid (Euprymna scolopes), which has symbiotic bacteria (Aliivibrio fischeri) that produce light to help the squid avoid nocturnal predators. This light shines through the squid's skin on its underside and is generated by a large and complex two-lobed light organ inside the squid's mantle cavity. From there, it escapes downwards, some of it travelling directly, some coming off a reflector at the top of the organ (dorsal side). Below there is a kind of iris, which has branches (diverticula) of its ink sac, with a lens below that; both the reflector and lens are derived from mesoderm. The squid controls light production by changing the shape of its iris or adjusting the strength of yellow filters on its underside, which presumably change the balance of wavelengths emitted. Light production shows a correlation with intensity of down-welling light, but it is about one third as bright; the squid can track repeated changes in brightness. Because the Hawaiian bobtail squid hides in sand during the day to avoid predators, it does not use counter-illumination during daylight hours. Predator distraction with ink Squid distract attacking predators by ejecting a cloud of ink, giving themselves an opportunity to escape. The ink gland and its associated ink sac empties into the rectum close to the anus, allowing the squid to rapidly discharge black ink into the mantle cavity and surrounding water. The ink is a suspension of melanin particles and quickly disperses to form a dark cloud that obscures the escape manoeuvres of the squid. Predatory fish may also be deterred by the alkaloid nature of the discharge which may interfere with their chemoreceptors. Nervous system and sense organs Cephalopods have the most highly developed nervous systems among invertebrates. Squids have a complex brain in the form of a nerve ring encircling the oesophagus, enclosed in a cartilaginous cranium. Paired cerebral ganglia above the oesophagus receive sensory information from the eyes and statocysts, and further ganglia below control the muscles of the mouth, foot, mantle and viscera. Giant axons up to in diameter convey nerve messages with great rapidity to the circular muscles of the mantle wall, allowing a synchronous, powerful contraction and maximum speed in the jet propulsion system. The paired eyes, on either side of the head, are housed in capsules fused to the cranium. Their structure is very similar to that of a fish eye, with a globular lens that has a depth of focus from to infinity. The image is focused by changing the position of the lens, as in a camera or telescope, rather than changing the shape of the lens, as in the human eye. Squid adjust to changes in light intensity by expanding and contracting the slit-shaped pupil. Deep sea squids in the family Histioteuthidae have eyes of two different types and orientation. The large left eye is tubular in shape and looks upwards, presumably searching for the silhouettes of animals higher in the water column. The normally-shaped right eye points forwards and downwards to detect prey. The statocysts are involved in maintaining balance and are analogous to the inner ear of fish. They are housed in cartilaginous capsules on either side of the cranium. They provide the squid with information on its body position in relation to gravity, its orientation, acceleration and rotation, and are able to perceive incoming vibrations. Without the statocysts, the squid cannot maintain equilibrium. Squid appear to have limited hearing, but the head and arms bear lines of hair-cells that are weakly sensitive to water movements and changes in pressure, and are analogous in function to the lateral line system of fish. Reproductive system The sexes are separate in squid, with a single gonad in the posterior part of the body. Fertilisation is external and usually takes place in the mantle cavity of the female. The male has a testis from which sperm pass into a single gonoduct where they are rolled together into a long bundle, or spermatophore. The gonoduct is elongated into a "penis" that extends into the mantle cavity and through which spermatophores are ejected. In shallow water species, the penis is short, and the spermatophore is removed from the mantle cavity by a tentacle of the male, which is specially adapted for the purpose and known as a hectocotylus, and placed inside the mantle cavity of the female during mating. The female has a large translucent ovary, situated towards the posterior of the visceral mass. From here, eggs travel along the gonocoel, where there are a pair of white nidamental glands, which lie anterior to the gills. Also present are red-spotted accessory nidamental glands containing symbiotic bacteria; both organs are associated with nutrient manufacture and forming shells for the eggs. The gonocoel enters the mantle cavity at the gonopore, and in some species, receptacles for storing spermatophores are located nearby, in the mantle wall. In shallow-water species of the continental shelf and epipelagic or mesopelagic zones, it is frequently one or both of arm pair IV of males that are modified into hectocotyli. However, most deep-sea squid lack hectocotyl arms and have longer penises; Ancistrocheiridae and Cranchiinae are exceptions. Giant squid of the genus Architeuthis are unusual in that they possess both a large penis and modified arm tips, although whether the latter are used for spermatophore transfer is uncertain. Penis elongation has been observed in the deep-water species Onykia ingens; when erect, the penis may be as long as the mantle, head, and arms combined. As such, deep-water squid have the greatest known penis length relative to body size of all mobile animals, second in the entire animal kingdom only to certain sessile barnacles. Digestive system Like all cephalopods, squids are predators and have complex digestive systems. The mouth is equipped with a sharp, horny beak mainly made of chitin and cross-linked proteins, which is used to kill and tear prey into manageable pieces. The beak is very robust, but does not contain minerals, unlike the teeth and jaws of many other organisms; the cross-linked proteins are histidine- and glycine-rich and give the beak a stiffness and hardness greater than most equivalent synthetic organic materials. The stomachs of captured whales often have indigestible squid beaks inside. The mouth contains the radula, the rough tongue common to all molluscs except bivalvia, which is equipped with multiple rows of teeth. In some species, toxic saliva helps to control large prey; when subdued, the food can be torn in pieces by the beak, moved to the oesophagus by the radula, and swallowed. The food bolus is moved along the gut by waves of muscular contractions (peristalsis). The long oesophagus leads to a muscular stomach roughly in the middle of the visceral mass. The digestive gland, which is equivalent to a vertebrate liver, diverticulates here, as does the pancreas, and both of these empty into the caecum, a pouch-shaped sac where most of the absorption of nutrients takes place. Indigestible food can be passed directly from the stomach to the rectum where it joins the flow from the caecum and is voided through the anus into the mantle cavity. Cephalopods are short-lived, and in mature squid, priority is given to reproduction; the female Onychoteuthis banksii for example, sheds its feeding tentacles on reaching maturity, and becomes flaccid and weak after spawning. Cardiovascular and excretory systems The squid mantle cavity is a seawater-filled sac containing three hearts and other organs supporting circulation, respiration, and excretion. Squid have a main systemic heart that pumps blood around the body as part of the general circulatory system, and two branchial hearts. The systemic heart consists of three chambers, a lower ventricle and two upper atria, all of which can contract to propel the blood. The branchial hearts pump blood specifically to the gills for oxygenation, before returning it to the systemic heart. The blood contains the copper-rich protein hemocyanin, which is used for oxygen transport at low ocean temperatures and low oxygen concentrations, and makes the oxygenated blood a deep, blue color. As systemic blood returns via two venae cavae to the branchial hearts, excretion of urine, carbon dioxide, and waste solutes occurs through outpockets (called nephridial appendages) in the venae cavae walls that enable gas exchange and excretion via the mantle cavity seawater. Buoyancy Unlike nautiloids and cuttlefish which have gas-filled chambers inside their shells which provide buoyancy, and octopuses which live near and rest on the seabed and do not require to be buoyant, many squid have a fluid-filled receptacle, equivalent to the swim bladder of a fish, in the coelom or connective tissue. This reservoir acts as a chemical buoyancy chamber, with the heavy metallic cations typical of seawater replaced by low molecular-weight ammonium ions, a product of excretion. The small difference in density provides a small contribution to buoyancy per unit volume, so the mechanism requires a large buoyancy chamber to be effective. Since the chamber is filled with liquid, it has the advantage over a swim bladder of not changing significantly in volume with pressure. Glass squids in the family Cranchiidae for example, have an enormous transparent coelom containing ammonium ions and occupying about two-thirds the volume of the animal, allowing it to float at the required depth. About half of the 28 families of squid use this mechanism to solve their buoyancy issues. The family Bathyteuthidae get their buoyancy from an oily substance found in their liver and around their mantle and head. Largest and smallest The majority of squid are no more than long, although the giant squid may reach . The smallest species are probably the benthic pygmy squids Idiosepius, which grow to a mantle length of , and have short bodies and stubby arms. In 1978, sharp, curved claws on the suction cups of squid tentacles cut up the rubber coating on the hull of the USS Stein. The size suggested the largest squid known at the time. In 2003, a large specimen of an abundant but poorly understood species, Mesonychoteuthis hamiltoni (the colossal squid), was discovered. This species may grow to in length, making it the largest invertebrate. In February 2007, a New Zealand fishing vessel caught the largest squid ever documented, weighing and measuring around off the coast of Antarctica. Dissection showed that the eyes, used to detect prey in the deep Southern Ocean, exceeded the size of footballs; these may be among the largest eyes ever to exist in the animal kingdom. Development The eggs of squid are large for a mollusc, containing a large amount of yolk to nourish the embryo as it develops directly, without an intervening veliger larval stage. The embryo grows as a disc of cells on top of the yolk. During the gastrulation stage, the margins of the disc grow to surround the yolk, forming a yolk sac, which eventually forms part of the animal's gut. The dorsal side of the disc grows upwards and forms the embryo, with a shell gland on its dorsal surface, gills, mantle and eyes. The arms and funnel develop as part of the foot on the ventral side of the disc. The arms later migrate upwards, coming to form a ring around the funnel and mouth. The yolk is gradually absorbed as the embryo grows. Some juvenile squid live higher in the water column than do adults. Squids tend to be short-lived; Loligo for example lives from one to three years according to species, typically dying soon after spawning. In a well-studied bioluminescent species, the Hawaiian bobtail squid, a special light organ in the squid's mantle is rapidly colonized with Aliivibrio fischeri bacteria within hours of hatching. This light-organ colonization requires this particular bacterial species for a symbiotic relationship; no colonization occurs in the absence of A. fischeri. Colonization occurs in a horizontal manner, such that the hosts acquires its bacterial partners from the environment. The symbiosis is obligate for the squid, but facultative for the bacteria. Once the bacteria enter the squid, they colonize interior epithelial cells in the light organ, living in crypts with complex microvilli protrusions. The bacteria also interact with hemocytes, macrophage-like blood cells that migrate between epithelial cells, but the mechanism and function of this process is not well understood. Bioluminescence reaches its highest levels during the early evening hours and bottoms out before dawn; this occurs because at the end of each day, the contents of the squid's crypts are expelled into the surrounding environment. Approximately 95% of the bacteria are voided each morning before the bacterial population builds up again by nightfall. Behaviour Locomotion Squid can move about in several different ways. Slow movement is achieved by a gentle undulation of the muscular lateral fins on either side of the trunk which drives the animal forward. A more common means of locomotion providing sustained movement is achieved using jetting, during which contraction of the muscular wall of the mantle cavity provides jet propulsion. Slow jetting is used for ordinary locomotion, and ventilation of the gills is achieved at the same time. The circular muscles in the mantle wall contract; this causes the inhalant valve to close, the exhalant valve to open and the mantle edge to lock tightly around the head. Water is forced out through the funnel which is pointed in the opposite direction to the required direction of travel. The inhalant phase is initiated by the relaxation of the circular muscles causes them to stretch, the connective tissue in the mantle wall recoils elastically, the mantle cavity expands causing the inhalant valve to open, the exhalant valve to close and water to flow into the cavity. This cycle of exhalation and inhalation is repeated to provide continuous locomotion. Fast jetting is an escape response. In this form of locomotion, radial muscles in the mantle wall are involved as well as circular ones, making it possible to hyper-inflate the mantle cavity with a larger volume of water than during slow jetting. On contraction, water flows out with great force, the funnel always being pointed anteriorly, and travel is backwards. During this means of locomotion, some squid exit the water in a similar way to flying fish, gliding through the air for up to , and occasionally ending up on the decks of ships. Feeding Squid are carnivores, and, with their strong arms and suckers, can overwhelm relatively large animals efficiently. Prey is identified by sight or by touch, grabbed by the tentacles which can be shot out with great rapidity, brought back to within reach of the arms, and held by the hooks and suckers on their surface. In some species, the squid's saliva contains toxins which act to subdue the prey. These are injected into its bloodstream when the prey is bitten, along with vasodilators and chemicals to stimulate the heart, and quickly circulate to all parts of its body. The deep sea squid Taningia danae has been filmed releasing blinding flashes of light from large photophores on its arms to illuminate and disorientate potential prey. Although squid can catch large prey, the mouth is relatively small, and the food must be cut into pieces by the chitinous beak with its powerful muscles before being swallowed. The radula is located in the buccal cavity and has multiple rows of tiny teeth that draw the food backwards and grind it in pieces. The deep sea squid Mastigoteuthis has the whole length of its whip-like tentacles covered with tiny suckers; it probably catches small organisms in the same way that flypaper traps flies. The tentacles of some bathypelagic squids bear photophores which may bring food within its reach by attracting prey. Squid are among the most intelligent invertebrates. For example, groups of Humboldt squid hunt cooperatively, spiralling up through the water at night and coordinating their vertical and horizontal movements while foraging. Reproduction Courtship in squid takes place in the open water and involves the male selecting a female, the female responding, and the transfer by the male of spermatophores to the female. In many instances, the male may display to identify himself to the female and drive off any potential competitors. Elaborate changes in body patterning take place in some species in both agonistic and courtship behaviour. The Caribbean reef squid (Sepioteuthis sepioidea), for example, employs a complex array of colour changes during courtship and social interactions and has a range of about 16 body patterns in its repertoire. The pair adopt a head-to-head position, and "jaw locking" may take place, in a similar manner to that adopted by some cichlid fish. The heterodactylus of the male is used to transfer the spermatophore and deposit it in the female's mantle cavity in the position appropriate for the species; this may be adjacent to the gonopore or in a seminal receptacle. The sperm may be used immediately or may be stored. As the eggs pass down the oviduct, they are wrapped in a gelatinous coating, before continuing to the mantle cavity, where they are fertilised. In Loligo, further coatings are added by the nidimental glands in the walls of the cavity and the eggs leave through a funnel formed by the arms. The female attaches them to the substrate in strings or groups, the coating layers swelling and hardening after contact with sea water. Loligo sometimes forms breeding aggregations which may create a "community pile" of egg strings. Some pelagic and deep sea squid do not attach their egg masses, which float freely. Ecology Squid mostly have an annual life cycle, growing fast and dying soon after spawning. The diet changes as they grow but mostly consists of large zooplankton and small nekton. In Antarctica for example, krill is the main constituent of the diet, with other food items being amphipods, other small crustaceans, and large arrow worms. Fish are also eaten, and some squid are cannibalistic. As well as occupying a key role in the food chain, squid are an important prey for predators including sharks, sea birds, seals and whales. Juvenile squid provide part of the diet for worms and small fish. When researchers studied the contents of the stomachs of elephant seals in South Georgia, they found 96% squid by weight. In a single day, a sperm whale can eat 700 to 800 squid, and a Risso's dolphin entangled in a net in the Mediterranean was found to have eaten angel clubhook squid, umbrella squid, reverse jewel squid and European flying squid, all identifiable from their indigestible beaks. Ornithoteuthis volatilis, a common squid from the tropical Indo-Pacific, is predated by yellowfin tuna, longnose lancetfish, common dolphinfish and swordfish, the tiger shark, the scalloped hammerhead shark and the smooth hammerhead shark. Sperm whales also hunt this species extensively as does the brown fur seal. In the Southern Ocean, penguins and wandering albatrosses are major predators of Gonatus antarcticus. Human uses In literature and art Giant squid have featured as monsters of the deep since classical times. Giant squid were described by Aristotle (4th century BC) in his History of Animals and Pliny the Elder (1st century AD) in his Natural History. The Gorgon of Greek mythology may have been inspired by squid or octopus, the animal itself representing the severed head of Medusa, the beak as the protruding tongue and fangs, and its tentacles as the snakes. The six-headed sea monster of the Odyssey, Scylla, may have had a similar origin. The Nordic legend of the kraken may also have derived from sightings of large cephalopods. In literature, H. G. Wells' short story "The Sea Raiders" featured a man-eating squid species Haploteuthis ferox. The science fiction writer Jules Verne told a tale of a kraken-like monster in his 1870 novel Twenty Thousand Leagues Under the Seas. As food Squid form a major food resource and are used in cuisines around the world, notably in Japan where it is eaten as ika sōmen, sliced into vermicelli-like strips; as sashimi; and as tempura. Three species of Loligo are used in large quantities: L. vulgaris in the Mediterranean (known as in Spanish, in Italian); L. forbesii in the Northeast Atlantic; and L. pealei on the American East Coast. Among the Ommastrephidae, Todarodes pacificus is the main commercial species, harvested in large quantities across the North Pacific in Canada, Japan and China. In English-speaking countries, squid as food is often called calamari, adopted from Italian into English in the 17th century. Squid are found abundantly in certain areas, and provide large catches for fisheries. The body can be stuffed whole, cut into flat pieces, or sliced into rings. The arms, tentacles, and ink are also edible; the only parts not eaten are the beak and gladius (pen). Squid is a good food source for zinc and manganese, and high in copper, selenium, vitamin B12, and riboflavin. Commercial fishing According to the FAO, the cephalopod catch for 2002 was . Of this, 2,189,206 tonnes, or 75.8 percent, was squid. The following table lists squid species fishery catches that exceeded in 2002. In biomimicry Prototype chromatophores that mimic the squid's adaptive camouflage have been made by Bristol University researchers using an electroactive dielectric elastomer, a flexible "smart" material that changes its colour and texture in response to electrical signals. The researchers state that their goal is to create an artificial skin that provides rapid active camouflage. The squid giant axon inspired Otto Schmitt to develop a comparator circuit with hysteresis now called the Schmitt trigger, replicating the axon's propagation of nerve impulses.	Biology and health sciences	Mollusks	null
38074	https://en.wikipedia.org/wiki/Kidney%20stone%20disease	Kidney stone disease	Kidney stone disease, also known as renal calculus disease, nephrolithiasis or urolithiasis, is a crystallopathy where a solid piece of material (renal calculus) develops in the urinary tract. Renal calculi typically form in the kidney and leave the body in the urine stream. A small calculus may pass without causing symptoms. If a stone grows to more than , it can cause blockage of the ureter, resulting in sharp and severe pain in the lower back that often radiates downward to the groin (renal colic). A calculus may also result in blood in the urine, vomiting, or painful urination. About half of people who have had a renal calculus are likely to have another within ten years. Most calculi form by a combination of genetics and environmental factors. Risk factors include high urine calcium levels, obesity, certain foods, some medications, calcium supplements, hyperparathyroidism, gout and not drinking enough fluids. Calculi form in the kidney when minerals in urine are at high concentration. The diagnosis is usually based on symptoms, urine testing, and medical imaging. Blood tests may also be useful. Calculi are typically classified by their location: nephrolithiasis (in the kidney), ureterolithiasis (in the ureter), cystolithiasis (in the bladder), or by what they are made of (calcium oxalate, uric acid, struvite, cystine). In those who have had renal calculi, drinking fluids is a way to prevent them. Drinking fluids such that more than two liters of urine are produced per day is recommended. If fluid intake alone is not effective to prevent renal calculi, the medications thiazide diuretic, citrate, or allopurinol may be suggested. Soft drinks containing phosphoric acid (typically colas) should be avoided. When a calculus causes no symptoms, no treatment is needed. For those with symptoms, pain control is usually the first measure, using medications such as nonsteroidal anti-inflammatory drugs or opioids. Larger calculi may be helped to pass with the medication tamsulosin or may require procedures such as extracorporeal shock wave lithotripsy, ureteroscopy, or percutaneous nephrolithotomy. Renal calculi have affected humans throughout history with a description of surgery to remove them dating from as early as 600 BC in ancient India by Sushruta. Between 1% and 15% of people globally are affected by renal calculi at some point in their lives. In 2015, 22.1 million cases occurred, resulting in about 16,100 deaths. They have become more common in the Western world since the 1970s. Generally, more men are affected than women. The prevalence and incidence of the disease rises worldwide and continues to be challenging for patients, physicians, and healthcare systems alike. In this context, epidemiological studies are striving to elucidate the worldwide changes in the patterns and the burden of the disease and identify modifiable risk factors that contribute to the development of renal calculi. Signs and symptoms The hallmark of a stone that obstructs the ureter or renal pelvis is excruciating, intermittent pain that radiates from the flank to the groin or to the inner thigh. This is due to the transfer of referred pain signals from the lower thoracic splanchnic nerves to the lumbar splanchnic nerves as the stone passes down from the kidney or proximal ureter to the distal ureter. This pain, known as renal colic, is often described as one of the strongest pain sensations known. Renal colic caused by kidney stones is commonly accompanied by urinary urgency, restlessness, hematuria, sweating, nausea, and vomiting. It typically comes in waves lasting 20 to 60 minutes caused by peristaltic contractions of the ureter as it attempts to expel the stone. The embryological link between the urinary tract, the genital system, and the gastrointestinal tract is the basis of the radiation of pain to the gonads, as well as the nausea and vomiting that are also common in urolithiasis. Postrenal azotemia and hydronephrosis can be observed following the obstruction of urine flow through one or both ureters. Pain in the lower-left quadrant can sometimes be confused with diverticulitis because the sigmoid colon overlaps the ureter, and the exact location of the pain may be difficult to isolate due to the proximity of these two structures. Risk factors Dehydration from low fluid intake is a factor in stone formation. Individuals living in warm climates are at higher risk due to increased fluid loss. Obesity, immobility, and sedentary lifestyles are other leading risk factors. High dietary intake of animal protein, sodium, sugars including honey, refined sugars, fructose and high fructose corn syrup, and excessive consumption of fruit juices may increase the risk of kidney stone formation due to increased uric acid excretion and elevated urinary oxalate levels (whereas tea, coffee, wine and beer may decrease the risk). Kidney stones can result from an underlying metabolic condition, such as distal renal tubular acidosis, Dent's disease, hyperparathyroidism, primary hyperoxaluria, or medullary sponge kidney. 3–20% of people who form kidney stones have medullary sponge kidney. Kidney stones are more common in people with Crohn's disease; Crohn's disease is associated with hyperoxaluria and malabsorption of magnesium. A person with recurrent kidney stones may be screened for such disorders. This is typically done with a 24-hour urine collection. The urine is analyzed for features that promote stone formation. Calcium oxalate Calcium is one component of the most common type of human kidney stones, calcium oxalate. Some studies suggest that people who take calcium or vitamin D as a dietary supplement have a higher risk of developing kidney stones. In the United States, kidney stone formation was used as an indicator of excess calcium intake by the Reference Daily Intake committee for calcium in adults. In the early 1990s, a study conducted for the Women's Health Initiative in the US found that postmenopausal women who consumed 1000 mg of supplemental calcium and 400 international units of vitamin D per day for seven years had a 17% higher risk of developing kidney stones than subjects taking a placebo. The Nurses' Health Study also showed an association between supplemental calcium intake and kidney stone formation. Unlike supplemental calcium, high intakes of dietary calcium do not appear to cause kidney stones and may actually protect against their development. This is perhaps related to the role of calcium in binding ingested oxalate in the gastrointestinal tract. As the amount of calcium intake decreases, the amount of oxalate available for absorption into the bloodstream increases; this oxalate is then excreted in greater amounts into the urine by the kidneys. In the urine, oxalate is a very strong promoter of calcium oxalate precipitation—about 15 times stronger than calcium. A 2004 study found that diets low in calcium are associated with a higher overall risk for kidney stone formation. For most individuals, other risk factors for kidney stones, such as high intakes of dietary oxalates and low fluid intake, play a greater role than calcium intake. Other electrolytes Calcium is not the only electrolyte that influences the formation of kidney stones. For example, by increasing urinary calcium excretion, high dietary sodium may increase the risk of stone formation. Drinking fluoridated tap water may increase the risk of kidney stone formation by a similar mechanism, though further epidemiologic studies are warranted to determine whether fluoride in drinking water is associated with an increased incidence of kidney stones. High dietary intake of potassium appears to reduce the risk of stone formation because potassium promotes the urinary excretion of citrate, an inhibitor of calcium crystal formation. Kidney stones are more likely to develop, and to grow larger, if a person has low dietary magnesium. Magnesium inhibits stone formation. Animal protein Diets in Western nations typically contain a large proportion of animal protein. Eating animal protein creates an acid load that increases urinary excretion of calcium and uric acid and reduced citrate. Urinary excretion of excess sulfurous amino acids (e.g., cysteine and methionine), uric acid, and other acidic metabolites from animal protein acidifies the urine, which promotes the formation of kidney stones. Low urinary-citrate excretion is also commonly found in those with a high dietary intake of animal protein, whereas vegetarians tend to have higher levels of citrate excretion. Low urinary citrate, too, promotes stone formation. Vitamins The evidence linking vitamin C supplements with an increased rate of kidney stones is inconclusive. The excess dietary intake of vitamin C might increase the risk of calcium-oxalate stone formation. The link between vitamin D intake and kidney stones is also tenuous. Excessive vitamin D supplementation may increase the risk of stone formation by increasing the intestinal absorption of calcium; correction of a deficiency does not. Pathophysiology Supersaturation of urine Kidney stones are primarily composed of calcium salts, with the most common being calcium oxalate (70-80%), followed by calcium phosphate and uric acid. When urine contains high concentrations of these ions, they can form crystals and eventually stones. The formation of kidney stones occurs in three main phases: nucleation (initial crystal formation) growth (expansion of single crystals) aggregation (clumping together of multiple crystals) When the urine becomes supersaturated (when the urine solvent contains more solutes than it can hold in solution) with one or more calculogenic (crystal-forming) substances, initial seed crystals may form through the process of nucleation. Heterogeneous nucleation (where there is a solid surface present on which a crystal can grow) proceeds more rapidly than homogeneous nucleation (where a crystal must grow in a liquid medium with no such surface), because it requires less energy. Adhering to cells on the surface of a renal papilla, a seed crystal can grow and aggregate into an organized mass. Depending on the chemical composition of the crystal, the stone-forming process may proceed more rapidly when the urine pH is unusually high or low. Supersaturation of the urine with respect to a calculogenic compound is pH-dependent. For example, at a pH of 7.0, the solubility of uric acid in urine is 158 mg/100 mL. Reducing the pH to 5.0 decreases the solubility of uric acid to less than 8 mg/100 mL. The formation of uric-acid stones requires a combination of hyperuricosuria (high urine uric-acid levels) and low urine pH; hyperuricosuria alone is not associated with uric-acid stone formation if the urine pH is alkaline. Supersaturation of the urine is a necessary, but not a sufficient, condition for the development of any urinary calculus. Supersaturation is likely the underlying cause of uric acid and cystine stones, but calcium-based stones (especially calcium oxalate stones) may have a more complex cause. Randall's plaque While supersaturation of urine may lead to crystalluria, it does not necessarily promote the formation of a kidney stone because the particle may not reach the sufficient size needed for renal attachment. On the other hand, Randall's plaques, which were first identified by Alexander Randall in 1937, are calcium phosphate deposits that form in the papillary interstitium and are thought to be the nidus required for stone development. In addition to Randall's plugs, which form in the Duct of Bellini, these structures can generate reactive oxygen species that further enhance stone formation. Pathogenic bacteria Some bacteria have roles in promoting stone formation. Specifically, urease-positive bacteria, such as Proteus mirabilis can produce the enzyme urease, which converts urea to ammonia and carbon dioxide. This increases the urinary pH and promotes struvite stone formation. Additionally, non-urease producing bacteria can provide bacterial components that promote calcium oxalate crystallization, though this mechanism is poorly understood. Inhibitors of stone formation Normal urine contains chelating agents, such as citrate, that inhibit the nucleation, growth, and aggregation of calcium-containing crystals. Other endogenous inhibitors include calgranulin (an S-100 calcium-binding protein), Tamm–Horsfall protein, glycosaminoglycans, uropontin (a form of osteopontin), nephrocalcin (an acidic glycoprotein), prothrombin F1 peptide, and bikunin (uronic acid-rich protein). The biochemical mechanisms of action of these substances have not yet been thoroughly elucidated. However, when these substances fall below their normal proportions, stones can form from an aggregation of crystals. Sufficient dietary intake of magnesium and citrate inhibits the formation of calcium oxalate and calcium phosphate stones; in addition, magnesium and citrate operate synergistically to inhibit kidney stones. The efficacy of magnesium in subduing stone formation and growth is dose-dependent. Hypocitraturia Hypocitraturia or low urinary-citrate excretion (variably defined as less than 320 mg/day) can be a contributing cause of kidney stones in up to 2/3 of cases. The protective role of citrate is linked to several mechanisms; citrate reduces urinary supersaturation of calcium salts by forming soluble complexes with calcium ions and by inhibiting crystal growth and aggregation. Therapy with potassium citrate is commonly prescribed in clinical practice to increase urinary citrate and to reduce stone formation rates. Alkali citrate is also used to increase urine citrate levels. It can be prescribed or found over-the-counter in pill, liquid or powder form. Diagnosis Diagnosis of kidney stones is made on the basis of information obtained from the history, physical examination, urinalysis, and radiographic studies. Clinical diagnosis is usually made on the basis of the location and severity of the pain, which is typically colicky in nature (comes and goes in spasmodic waves). Pain in the back occurs when calculi produce an obstruction in the kidney. Physical examination may reveal fever and tenderness at the costovertebral angle on the affected side. Imaging studies Calcium-containing stones are relatively radiodense (opaque to X-rays), and they can often be detected by a traditional radiography of the abdomen that includes the kidneys, ureters, and bladder (KUB film). KUB radiography, although useful in monitoring size of stone or passage of stone in stone formers, might not be useful in the acute setting due to low sensitivity. Some 60% of all renal stones are radiopaque. In general, calcium phosphate stones have the greatest density, followed by calcium oxalate and magnesium ammonium phosphate stones. Cystine calculi are only faintly radiodense, while uric acid stones are usually entirely radiolucent. In people with a history of stones, those who are less than 50 years of age and are presenting with the symptoms of stones without any concerning signs do not require helical CT scan imaging. A computed tomography (CT) scan is also not typically recommended in children. Otherwise a noncontrast helical CT scan with sections is the diagnostic method to use to detect kidney stones and confirm the diagnosis of kidney stone disease. Near all stones are detectable on CT scans with the exception of those composed of certain drug residues in the urine, such as from indinavir. Where a CT scan is unavailable, an intravenous pyelogram may be performed to help confirm the diagnosis of urolithiasis. This involves intravenous injection of a contrast agent followed by a KUB film. Uroliths present in the kidneys, ureters, or bladder may be better defined by the use of this contrast agent. Stones can also be detected by a retrograde pyelogram, where a similar contrast agent is injected directly into the distal ostium of the ureter (where the ureter terminates as it enters the bladder). Renal ultrasonography can sometimes be useful, because it gives details about the presence of hydronephrosis, suggesting that the stone is blocking the outflow of urine. Radiolucent stones, which do not appear on KUB, may show up on ultrasound imaging studies. Other advantages of renal ultrasonography include its low cost and absence of radiation exposure. Ultrasound imaging is useful for detecting stones in situations where X-rays or CT scans are discouraged, such as in children or pregnant women. Despite these advantages, renal ultrasonography in 2009 was not considered a substitute for noncontrast helical CT scan in the initial diagnostic evaluation of urolithiasis. The main reason for this is that, compared with CT, renal ultrasonography more often fails to detect small stones (especially ureteral stones) and other serious disorders that could be causing the symptoms. On the contrary, a 2014 study suggested that ultrasonography should be used as the initial diagnostic imaging test, with further imaging studies be performed at the discretion of the physician on the basis of clinical judgment, and using ultrasonography rather than CT as an initial diagnostic test results in less radiation exposure and equally good outcome. Laboratory examination Laboratory investigations typically carried out include: microscopic examination of the urine, which may show red blood cells, bacteria, leukocytes, urinary casts, and crystals; urine culture to identify any infecting organisms present in the urinary tract and sensitivity to determine the susceptibility of these organisms to specific antibiotics; complete blood count, looking for neutrophilia (increased neutrophil granulocyte count) suggestive of bacterial infection, as seen in the setting of struvite stones; renal function tests to look for abnormally high blood calcium levels (hypercalcemia); 24 hour urine collection to measure total daily urinary volume, magnesium, sodium, uric acid, calcium, citrate, oxalate, and phosphate; collection of stones (by urinating through a StoneScreen kidney stone collection cup or a simple tea strainer) is useful. Chemical analysis of collected stones can establish their composition, which in turn can help to guide future preventive and therapeutic management. Composition Calcium-containing stones By far, the most common type of kidney stones worldwide contains calcium. For example, calcium-containing stones represent about 80% of all cases in the United States; these typically contain calcium oxalate either alone or in combination with calcium phosphate in the form of apatite or brushite. Factors that promote the precipitation of oxalate crystals in the urine, such as primary hyperoxaluria, are associated with the development of calcium oxalate stones. The formation of calcium phosphate stones is associated with conditions such as hyperparathyroidism and renal tubular acidosis. Oxaluria is increased in patients with certain gastrointestinal disorders including inflammatory bowel disease such as Crohn's disease or in patients who have undergone resection of the small bowel or small-bowel bypass procedures. Oxaluria is also increased in patients who consume increased amounts of oxalate (found in vegetables and nuts). Primary hyperoxaluria is a rare autosomal recessive condition that usually presents in childhood. Calcium oxalate crystals can come in two varieties. Calcium oxalate monohydrate can appear as 'dumbbells' or as long ovals that resemble the individual posts in a picket fence. Calcium oxalate dihydrate have a tetragonal "envelope" appearance. Struvite stones About 10–15% of urinary calculi are composed of struvite (hexa-hydrated ammonium magnesium phosphate, NH4MgPO4·6H2O). Struvite stones (also known as "infection stones," urease, or triple-phosphate stones) form most often in the presence of infection by urea-splitting bacteria. Using the enzyme urease, these organisms metabolize urea into ammonia and carbon dioxide. This alkalinizes the urine, resulting in favorable conditions for the formation of struvite stones. Proteus mirabilis, Proteus vulgaris, and Morganella morganii are the most common organisms isolated; less common organisms include Ureaplasma urealyticum and some species of Providencia, Klebsiella, Serratia, and Enterobacter. These infection stones are commonly observed in people who have factors that predispose them to urinary tract infections, such as those with spinal cord injury and other forms of neurogenic bladder, ileal conduit urinary diversion, vesicoureteral reflux, and obstructive uropathies. They are also commonly seen in people with underlying metabolic disorders, such as idiopathic hypercalciuria, hyperparathyroidism, and gout. Infection stones can grow rapidly, forming large calyceal staghorn (antler-shaped) calculi requiring invasive surgery such as percutaneous nephrolithotomy for definitive treatment. Struvite stones (triple-phosphate/magnesium ammonium phosphate) have a 'coffin lid' morphology by microscopy. Uric acid stones About 5–10% of all stones are formed from uric acid. People with certain metabolic abnormalities, including obesity, may produce uric acid stones. They also may form in association with conditions that cause hyperuricosuria (an excessive amount of uric acid in the urine) with or without hyperuricemia (an excessive amount of uric acid in the serum). They may also form in association with disorders of acid/base metabolism where the urine is excessively acidic (low pH), resulting in precipitation of uric acid crystals. A diagnosis of uric acid urolithiasis is supported by the presence of a radiolucent stone in the face of persistent urine acidity, in conjunction with the finding of uric acid crystals in fresh urine samples. As noted above (section on calcium oxalate stones), people with inflammatory bowel disease (Crohn's disease, ulcerative colitis) tend to have hyperoxaluria and form oxalate stones. They also have a tendency to form urate stones. Urate stones are especially common after colon resection. Uric acid stones appear as pleomorphic crystals, usually diamond-shaped. They may also look like squares or rods which are polarizable. Other types People with certain rare inborn errors of metabolism have a propensity to accumulate crystal-forming substances in their urine. For example, those with cystinuria, cystinosis, and Fanconi syndrome may form stones composed of cystine. Cystine stone formation can be treated with urine alkalinization and dietary protein restriction. People affected by xanthinuria often produce stones composed of xanthine. People affected by adenine phosphoribosyltransferase deficiency may produce 2,8-dihydroxyadenine stones, alkaptonurics produce homogentisic acid stones, and iminoglycinurics produce stones of glycine, proline, and hydroxyproline. Urolithiasis has also been noted to occur in the setting of therapeutic drug use, with crystals of drug forming within the renal tract in some people currently being treated with agents such as indinavir, sulfadiazine, and triamterene. Location Urolithiasis refers to stones originating anywhere in the urinary system, including the kidneys and bladder. Nephrolithiasis refers to the presence of such stones in the kidneys. Calyceal calculi are aggregations in either the minor or major calyx, parts of the kidney that pass urine into the ureter (the tube connecting the kidneys to the urinary bladder). The condition is called ureterolithiasis when a calculus is located in the ureter. Stones may also form or pass into the bladder, a condition referred to as bladder stones. Size Stones less than in diameter pass spontaneously in up to 98% of cases, while those measuring in diameter pass spontaneously in less than 53% of cases. Stones that are large enough to fill out the renal calyces are called staghorn stones and are composed of struvite in a vast majority of cases, which forms only in the presence of urease-forming bacteria. Other forms that can possibly grow to become staghorn stones are those composed of cystine, calcium oxalate monohydrate, and uric acid. Prevention Preventative measures depend on the type of stones. In those with calcium stones, drinking plenty of fluids, thiazide diuretics and citrate are effective as is allopurinol in those with high uric acid levels in urine. Dietary measures Specific therapy should be tailored to the type of stones involved. Diet can have an effect on the development of kidney stones. Preventive strategies include some combination of dietary modifications and medications with the goal of reducing the excretory load of calculogenic compounds on the kidneys. Dietary recommendations to minimize the formation of kidney stones include: increasing total fluid intake to achieve more than two liters per day of urine output; limiting cola, including sugar-sweetened soft drinks; to less than one liter per week. limiting animal protein intake to no more than two meals daily (an association between animal protein and recurrence of kidney stones has been shown in men); increasing citrate, including from lemon and lime juice; citric acid in its natural form, such as from citrus fruits, "prevents small stones from becoming 'problem stones' by coating them and preventing other material from attaching and building onto the stones"; citrate inhibits the formation of kidney stones on all phasesnucleation, growth and aggregationby raising the limit at which oxalate remain stable, slowing oxalate crystal growth, and notably, reducing crystal aggregation within the kidney tubules; increase alkaline load by consuming more fruits and vegetables (because uric acid crystals form in acidic environment); reducing sodium intake is associated with a reduction in urine calcium excretion. Maintenance of dilute urine by means of vigorous fluid therapy is beneficial in all forms of kidney stones, so increasing urine volume is a key principle for the prevention of kidney stones. Fluid intake should be sufficient to maintain a urine output of at least per day. A high fluid intake may reduce the likelihood of kidney stone recurrence or may increase the time between stone development without unwanted effects. Calcium binds with available oxalate in the gastrointestinal tract, thereby preventing its absorption into the bloodstream. Reducing oxalate absorption decreases kidney stone risk in susceptible people. Because of this, some doctors recommend increasing dairy intake so that its calcium content will serve as an oxalate binder. Taking calcium citrate tablets during or after meals containing high oxalate foods may be useful if dietary calcium cannot be increased by other means as in those with lactose intolerance. The preferred calcium supplement for people at risk of stone formation is calcium citrate, as opposed to calcium carbonate, because it helps to increase urinary citrate excretion. Aside from vigorous oral hydration and eating more dietary calcium, other prevention strategies include avoidance of higher doses of supplemental (since ascorbate is metabolized to oxalate) and restriction of oxalate-rich foods such as leaf vegetables, rhubarb, soy products and chocolate. However, no randomized, controlled trial of oxalate restriction has been performed to test the hypothesis that oxalate restriction reduces stone formation. Some evidence indicates magnesium intake decreases the risk of symptomatic kidney stones. Urine alkalinization The mainstay for medical management of uric acid stones is alkalinization (increasing the pH) of the urine. Uric acid stones are among the few types amenable to dissolution therapy, referred to as chemolysis. Chemolysis is usually achieved through the use of oral medications, although in some cases, intravenous agents or even instillation of certain irrigating agents directly onto the stone can be performed, using antegrade nephrostomy or retrograde ureteral catheters. Acetazolamide is a medication that alkalinizes the urine. In addition to acetazolamide or as an alternative, certain dietary supplements are available that produce a similar alkalinization of the urine. These include alkali citrate, sodium bicarbonate, potassium citrate, magnesium citrate, and bicitrate (a combination of citric acid monohydrate and sodium citrate dihydrate). Aside from alkalinization of the urine, these supplements have the added advantage of increasing the urinary citrate level, which helps to reduce the aggregation of calcium oxalate stones. Increasing the urine pH to around 6.5 provides optimal conditions for dissolution of uric acid stones. Increasing the urine pH to a value higher than 7.0 may increase the risk of calcium phosphate stone formation, though this concept is controversial since citrate does inhibit calcium phosphate crystallization. Testing the urine periodically with nitrazine paper can help to ensure the urine pH remains in this optimal range. Using this approach, stone dissolution rate can be expected to be around of stone radius per month. Slaked lime Calcium hydroxide decreases urinary calcium when combined with food rich in oxalic acid such as green leafy vegetables. Diuretics One of the recognized medical therapies for prevention of stones is the thiazide and thiazide-like diuretics, such as chlorthalidone or indapamide. These drugs inhibit the formation of calcium-containing stones by reducing urinary calcium excretion. Sodium restriction is necessary for clinical effect of thiazides, as sodium excess promotes calcium excretion. Thiazides work best for renal leak hypercalciuria (high urine calcium levels), a condition in which high urinary calcium levels are caused by a primary kidney defect. Thiazides are useful for treating absorptive hypercalciuria, a condition in which high urinary calcium is a result of excess absorption from the gastrointestinal tract. Allopurinol For people with hyperuricosuria and calcium stones, allopurinol is one of the few treatments that have been shown to reduce kidney stone recurrences. Allopurinol interferes with the production of uric acid in the liver. The drug is also used in people with gout or hyperuricemia (high serum uric acid levels). Dosage is adjusted to maintain a reduced urinary excretion of uric acid. Serum uric acid level at or below 6 mg/100 mL is often a therapeutic goal. Hyperuricemia is not necessary for the formation of uric acid stones; hyperuricosuria can occur in the presence of normal or even low serum uric acid. Some practitioners advocate adding allopurinol only in people in whom hyperuricosuria and hyperuricemia persist, despite the use of a urine-alkalinizing agent such as sodium bicarbonate or potassium citrate. Treatment Stone size influences the rate of spontaneous stone passage. For example, up to 98% of small stones (less than in diameter) may pass spontaneously through urination within four weeks of the onset of symptoms, but for larger stones ( in diameter), the rate of spontaneous passage decreases to less than 53%. Initial stone location also influences the likelihood of spontaneous stone passage. Rates increase from 48% for stones located in the proximal ureter to 79% for stones located at the vesicoureteric junction, regardless of stone size. Assuming no high-grade obstruction or associated infection is found in the urinary tract, and symptoms are relatively mild, various nonsurgical measures can be used to encourage the passage of a stone. Repeat stone formers benefit from more intense management, including proper fluid intake and use of certain medications, as well as careful monitoring. Pain management Management of pain often requires intravenous administration of NSAIDs or opioids. NSAIDs appear somewhat better than opioids or paracetamol in those with normal kidney function. Medications by mouth are often effective for less severe discomfort. The use of antispasmodics does not have further benefit. Medical expulsive therapy The use of medications to speed the spontaneous passage of stones in the ureter is referred to as medical expulsive therapy. Several agents, including alpha adrenergic blockers (such as tamsulosin) and calcium channel blockers (such as nifedipine), may be effective. Alpha-blockers likely result in more people passing their stones, and they may pass their stones in a shorter time. People taking alpha-blockers may also use less pain medication and may not need to visit the hospital. Alpha-blockers appear to be more effective for larger stones (over 5 mm in size) than smaller stones. However, use of alpha-blockers may be associated with a slight increase in serious, unwanted effects from this medication. A combination of tamsulosin and a corticosteroid may be better than tamsulosin alone. These treatments also appear to be useful in addition to lithotripsy. Lithotripsy Extracorporeal shock wave lithotripsy (ESWL) is a noninvasive technique for the removal of kidney stones. Most ESWL is carried out when the stone is present near the renal pelvis. ESWL involves the use of a lithotriptor machine to deliver externally applied, focused, high-intensity pulses of ultrasonic energy to cause fragmentation of a stone over a period of around 30–60 minutes. Following its introduction in the United States in February 1984, ESWL was rapidly and widely accepted as a treatment alternative for renal and ureteral stones. It is currently used in the treatment of uncomplicated stones located in the kidney and upper ureter, provided the aggregate stone burden (stone size and number) is less than and the anatomy of the involved kidney is normal. For a stone greater than , ESWL may not help break the stone in one treatment; instead, two or three treatments may be needed. Some 80-85% of simple renal calculi can be effectively treated with ESWL. A number of factors can influence its efficacy, including chemical composition of the stone, presence of anomalous renal anatomy and the specific location of the stone within the kidney, presence of hydronephrosis, body mass index, and distance of the stone from the surface of the skin. Common adverse effects of ESWL include acute trauma, such as bruising at the site of shock administration, and damage to blood vessels of the kidney. In fact, the vast majority of people who are treated with a typical dose of shock waves using currently accepted treatment settings are likely to experience some degree of acute kidney injury. ESWL-induced acute kidney injury is dose-dependent (increases with the total number of shock waves administered and with the power setting of the lithotriptor) and can be severe, including internal bleeding and subcapsular hematomas. On rare occasions, such cases may require blood transfusion and even lead to acute kidney failure. Hematoma rates may be related to the type of lithotriptor used; hematoma rates of less than 1% and up to 13% have been reported for different lithotriptor machines. Recent studies show reduced acute tissue injury when the treatment protocol includes a brief pause following the initiation of treatment, and both improved stone breakage and a reduction in injury when ESWL is carried out at slow shock wave rate. In addition to the aforementioned potential for acute kidney injury, animal studies suggest these acute injuries may progress to scar formation, resulting in loss of functional renal volume. Recent prospective studies also indicate elderly people are at increased risk of developing new-onset hypertension following ESWL. In addition, a retrospective case-control study published by researchers from the Mayo Clinic in 2006 has found an increased risk of developing diabetes mellitus and hypertension in people who had undergone ESWL, compared with age and gender-matched people who had undergone nonsurgical treatment. Whether or not acute trauma progresses to long-term effects probably depends on multiple factors that include the shock wave dose (i.e., the number of shock waves delivered, rate of delivery, power setting, acoustic characteristics of the particular lithotriptor, and frequency of retreatment), as well as certain intrinsic predisposing pathophysiologic risk factors. To address these concerns, the American Urological Association established the Shock Wave Lithotripsy Task Force to provide an expert opinion on the safety and risk-benefit ratio of ESWL. The task force published a white paper outlining their conclusions in 2009. They concluded the risk-benefit ratio remains favorable for many people. The advantages of ESWL include its noninvasive nature, the fact that it is technically easy to treat most upper urinary tract calculi, and that, at least acutely, it is a well-tolerated, low-morbidity treatment for the vast majority of people. However, they recommended slowing the shock wave firing rate from 120 pulses per minute to 60 pulses per minute to reduce the risk of renal injury and increase the degree of stone fragmentation. Alpha-blockers are sometimes prescribed after shock wave lithotripsy to help the pieces of the stone leave the person's body. By relaxing muscles and helping to keep blood vessels open, alpha blockers may relax the ureter muscles to allow the kidney stone fragments to pass. When compared to usual care or placebo treatment, alpha blockers may lead to faster clearing of stones, a reduced need for extra treatment and fewer unwanted effects. They may also clear kidney stones in more adults than the standard shock wave lithotripsy procedure. The unwanted effects associated with alpha blockers are hospital emergency visits and return to hospital for stone-related issues, but these effects were more common in adults who did not receive alpha-blockers as a part of their treatment. Surgery Most stones under pass spontaneously. Prompt surgery may, nonetheless, be required in persons with only one working kidney, bilateral obstructing stones, a urinary tract infection and thus, it is presumed, an infected kidney, or intractable pain. Beginning in the mid-1980s, less invasive treatments such as extracorporeal shock wave lithotripsy, ureteroscopy, and percutaneous nephrolithotomy began to replace open surgery as the modalities of choice for the surgical management of urolithiasis. More recently, flexible ureteroscopy has been adapted to facilitate retrograde nephrostomy creation for percutaneous nephrolithotomy. This approach is still under investigation, though early results are favorable. Percutaneous nephrolithotomy or, rarely, anatrophic nephrolithotomy, is the treatment of choice for large or complicated stones (such as calyceal staghorn calculi) or stones that cannot be extracted using less invasive procedures. Ureteroscopic surgery Ureteroscopy has become increasingly popular as flexible and rigid fiberoptic ureteroscopes have become smaller. One ureteroscopic technique involves the placement of a ureteral stent (a small tube extending from the bladder, up the ureter and into the kidney) to provide immediate relief of an obstructed kidney. Stent placement can be useful for saving a kidney at risk for postrenal acute kidney failure due to the increased hydrostatic pressure, swelling and infection (pyelonephritis and pyonephrosis) caused by an obstructing stone. Ureteral stents vary in length from and most have a shape commonly referred to as a "double-J" or "double pigtail", because of the curl at both ends. They are designed to allow urine to flow past an obstruction in the ureter. They may be retained in the ureter for days to weeks as infections resolve and as stones are dissolved or fragmented by ESWL or by some other treatment. The stents dilate the ureters, which can facilitate instrumentation, and they also provide a clear landmark to aid in the visualization of the ureters and any associated stones on radiographic examinations. The presence of indwelling ureteral stents may cause minimal to moderate discomfort, frequency or urgency incontinence, and infection, which in general resolves on removal. Most ureteral stents can be removed cystoscopically during an office visit under topical anesthesia after resolution of urolithiasis. Research is currently uncertain if placing a temporary stent during ureteroscopy leads to different outcomes than not placing a stent in terms of number of hospital visits for post operative problems, short or long term pain, need for narcotic pain medication, risk of UTI, need for a repeat procedure or narrowing of the ureter from scarring. More definitive ureteroscopic techniques for stone extraction (rather than simply bypassing the obstruction) include basket extraction and ultrasound ureterolithotripsy. Laser lithotripsy is another technique, which involves the use of a holmium:yttrium aluminium garnet (Ho:YAG) laser to fragment stones in the bladder, ureters, and kidneys. Ureteroscopic techniques are generally more effective than ESWL for treating stones located in the lower ureter, with success rates of 93–100% using Ho:YAG laser lithotripsy. Although ESWL has been traditionally preferred by many practitioners for treating stones located in the upper ureter, more recent experience suggests ureteroscopic techniques offer distinct advantages in the treatment of upper ureteral stones. Specifically, the overall success rate is higher, fewer repeat interventions and postoperative visits are needed, and treatment costs are lower after ureteroscopic treatment when compared with ESWL. These advantages are especially apparent with stones greater than in diameter. However, because ureteroscopy of the upper ureter is much more challenging than ESWL, many urologists still prefer to use ESWL as a first-line treatment for stones of less than 10 mm, and ureteroscopy for those greater than 10 mm in diameter. Ureteroscopy is the preferred treatment in pregnant and morbidly obese people, as well as those with bleeding disorders. Epidemiology Kidney stones affect all geographical, cultural, and racial groups. The lifetime risk is about 10-15% in the developed world, but can be as high as 20-25% in the Middle East. The increased risk of dehydration in hot climates, coupled with a diet 50% lower in calcium and 250% higher in oxalates compared to Western diets, accounts for the higher net risk in the Middle East. In the Middle East, uric acid stones are more common than calcium-containing stones. The number of deaths due to kidney stones is estimated at 19,000 per year being fairly consistent between 1990 and 2010. In North America and Europe, the annual number of new cases per year of kidney stones is roughly 0.5%. In the United States, the frequency in the population of urolithiasis has increased from 3.2% to 5.2% from the mid-1970s to the mid-1990s. In the United States, about 9% of the population has had a kidney stone. The total cost for treating urolithiasis was US$2 billion in 2003. About 65–80% of those with kidney stones are men; most stones in women are due to either metabolic defects (such as cystinuria) or infections in the case of struvite stones. Urinary tract calculi disorders are more common in men than in women. Men most commonly experience their first episode between 30 and 40 years of age, whereas for women, the age at first presentation is somewhat later. The age of onset shows a bimodal distribution in women, with episodes peaking at 35 and 55 years. Recurrence rates are estimated at 50% over a 10-year and 75% over 20-year period, with some people experiencing ten or more episodes over the course of a lifetime. A 2010 review concluded that rates of disease are increasing. History The existence of kidney stones was first recorded thousands of years ago, with various explanations given; Joseph Glanville's Saducismus Triumphatus, for example, gives a detailed description of Abraham Mechelburg's voiding of small stones through his penis' virga, attributing the issue to witchcraft. In 1901, a stone discovered in the pelvis of an ancient Egyptian mummy was dated to 4,800 BC. Medical texts from ancient Mesopotamia, India, China, Persia, Greece, and Rome all mentioned calculous disease. Part of the Hippocratic Oath suggests there were practicing surgeons in ancient Greece to whom physicians were to defer for lithotomies, or the surgical removal of stones. The Roman medical treatise De Medicina by Aulus Cornelius Celsus contained a description of lithotomy, and this work served as the basis for this procedure until the 18th century. Examples of people who had kidney stone disease include Napoleon I, Epicurus, Napoleon III, Peter the Great, Louis XIV, George IV, Oliver Cromwell, Lyndon B. Johnson, Benjamin Franklin, Michel de Montaigne, Francis Bacon, Isaac Newton, Samuel Pepys, William Harvey, Herman Boerhaave, and Antonio Scarpa. New techniques in lithotomy began to emerge starting in 1520, but the operation remained risky. After Henry Jacob Bigelow popularized the technique of litholapaxy in 1878, the mortality rate dropped from about 24% to 2.4%. However, other treatment techniques continued to produce a high level of mortality, especially among inexperienced urologists. In 1980, Dornier MedTech introduced extracorporeal shock wave lithotripsy for breaking up stones via acoustical pulses, and this technique has since come into widespread use. Etymology The term renal calculus is from the Latin rēnēs, meaning "kidneys", and calculus, meaning "pebble". Lithiasis (stone formation) in the kidneys is called nephrolithiasis (), from nephro-, meaning kidney, + -lith, meaning stone, and -iasis, meaning disorder. A distinction between nephrolithiasis and urolithiasis can be made because not all urinary stones (uroliths) form in the kidney; they can also form in the bladder. But the distinction is often clinically irrelevant (with similar disease process and treatment either way) and the words are thus often used loosely as synonyms. Children Although kidney stones do not often occur in children, the incidence is increasing. These stones are in the kidney in two thirds of reported cases, and in the ureter in the remaining cases. Older children are at greater risk independent of whether or not they are male or female. As with adults, most pediatric kidney stones are predominantly composed of calcium oxalate; struvite and calcium phosphate stones are less common. Calcium oxalate stones in children are associated with high amounts of calcium, oxalate, and magnesium in acidic urine. Treatment of kidney stones in children is similar to treatments for adults, including shock wave lithotripsy, medication, and treatment using scope through the bladder, kidney or skin. Of these treatments, research is uncertain if shock waves are more effective than medication or a scope through the bladder, but it is likely less successful than a scope through skin into the kidney. When going in with a scope through the kidney, a regular and a mini-sized scope likely have similar success rates of stone removal. Alpha-blockers, a type of medication, may increase the successful removal of kidney stones when compared with a placebo and without ibuprofen. Research Metabolic syndrome and its associated diseases of obesity and diabetes as general risk factors for kidney stone disease are under research to determine if urinary excretion of calcium, oxalate and urate are higher than in people with normal weight or underweight, and if diet and physical activity have roles. Dietary, fluid intake, and lifestyle factors remain major topics for research on prevention of kidney stones, as of 2017. Gut microbiota The gut microbiota has been explored as a contributing factor for stone disease, indicating that some bacteria may be different in people forming kidney stones. One bacterium, Oxalobacter formigenes, is potentially beneficial for mitigating calcium oxalate stones because of its ability to metabolize oxalate as its sole carbon source, but 2018 research suggests that it is instead part of a network of oxalate degrading bacteria. Additionally, one study found that oral antibiotic use, which alters the gut microbiota, can increase the odds of a person developing a kidney stone. In other animals Among ruminants, uroliths more commonly cause problems in males than in females; the sigmoid flexure of the ruminant male urinary tract is more likely to obstruct passage. Early-castrated males are at greater risk, because of lesser urethral diameter. Low Ca:P intake ratio is conducive to phosphatic (e.g. struvite) urolith formation. Incidence among wether lambs can be minimized by maintaining a dietary Ca:P intake ratio of 2:1. Alkaline (higher) pH favors formation of carbonate and phosphate calculi. For domestic ruminants, dietary cation: anion balance is sometimes adjusted to assure a slightly acidic urine pH, for prevention of calculus formation. Differing generalizations regarding effects of pH on formation of silicate uroliths may be found. In this connection, it may be noted that under some circumstances, calcium carbonate accompanies silica in siliceous uroliths. Pelleted feeds may be conducive to formation of phosphate uroliths, because of increased urinary phosphorus excretion. This is attributable to lower saliva production where pelleted rations containing finely ground constituents are fed. With less blood phosphate partitioned into saliva, more tends to be excreted in urine. (Most saliva phosphate is fecally excreted.) Oxalate uroliths can occur in ruminants, although such problems from oxalate ingestion may be relatively uncommon. Ruminant urolithiasis associated with oxalate ingestion has been reported. However, no renal tubular damage or visible deposition of calcium oxalate crystals in kidneys was found in yearling wether sheep fed diets containing soluble oxalate at 6.5 percent of dietary dry matter for about 100 days. Conditions limiting water intake can result in stone formation. Various surgical interventions, e.g. amputation of the urethral process at its base near the glans penis in male ruminants, perineal urethrostomy, or tube cystostomy may be considered for relief of obstructive urolithiasis.	Biology and health sciences	Specific diseases	Health
38128	https://en.wikipedia.org/wiki/Cauchy%E2%80%93Schwarz%20inequality	Cauchy–Schwarz inequality	The Cauchy–Schwarz inequality (also called Cauchy–Bunyakovsky–Schwarz inequality) is an upper bound on the inner product between two vectors in an inner product space in terms of the product of the vector norms. It is considered one of the most important and widely used inequalities in mathematics. Inner products of vectors can describe finite sums (via finite-dimensional vector spaces), infinite series (via vectors in sequence spaces), and integrals (via vectors in Hilbert spaces). The inequality for sums was published by . The corresponding inequality for integrals was published by and . Schwarz gave the modern proof of the integral version. Statement of the inequality The Cauchy–Schwarz inequality states that for all vectors and of an inner product space where is the inner product. Examples of inner products include the real and complex dot product; see the examples in inner product. Every inner product gives rise to a Euclidean norm, called the or , where the norm of a vector is denoted and defined by where is always a non-negative real number (even if the inner product is complex-valued). By taking the square root of both sides of the above inequality, the Cauchy–Schwarz inequality can be written in its more familiar form in terms of the norm: Moreover, the two sides are equal if and only if and are linearly dependent. Special cases Sedrakyan's lemma – positive real numbers Sedrakyan's inequality, also known as Bergström's inequality, Engel's form, Titu's lemma (or the T2 lemma), states that for real numbers and positive real numbers : or, using summation notation, It is a direct consequence of the Cauchy–Schwarz inequality, obtained by using the dot product on upon substituting and . This form is especially helpful when the inequality involves fractions where the numerator is a perfect square. - The plane The real vector space denotes the 2-dimensional plane. It is also the 2-dimensional Euclidean space where the inner product is the dot product. If and then the Cauchy–Schwarz inequality becomes: where is the angle between and . The form above is perhaps the easiest in which to understand the inequality, since the square of the cosine can be at most 1, which occurs when the vectors are in the same or opposite directions. It can also be restated in terms of the vector coordinates , , , and as where equality holds if and only if the vector is in the same or opposite direction as the vector , or if one of them is the zero vector. : n-dimensional Euclidean space In Euclidean space with the standard inner product, which is the dot product, the Cauchy–Schwarz inequality becomes: The Cauchy–Schwarz inequality can be proved using only elementary algebra in this case by observing that the difference of the right and the left hand side is or by considering the following quadratic polynomial in Since the latter polynomial is nonnegative, it has at most one real root, hence its discriminant is less than or equal to zero. That is, : n-dimensional complex space If with and (where and ) and if the inner product on the vector space is the canonical complex inner product (defined by where the bar notation is used for complex conjugation), then the inequality may be restated more explicitly as follows: That is, For the inner product space of square-integrable complex-valued functions, the following inequality holds. The Hölder inequality is a generalization of this. Applications Analysis In any inner product space, the triangle inequality is a consequence of the Cauchy–Schwarz inequality, as is now shown: Taking square roots gives the triangle inequality: The Cauchy–Schwarz inequality is used to prove that the inner product is a continuous function with respect to the topology induced by the inner product itself. Geometry The Cauchy–Schwarz inequality allows one to extend the notion of "angle between two vectors" to any real inner-product space by defining: The Cauchy–Schwarz inequality proves that this definition is sensible, by showing that the right-hand side lies in the interval and justifies the notion that (real) Hilbert spaces are simply generalizations of the Euclidean space. It can also be used to define an angle in complex inner-product spaces, by taking the absolute value or the real part of the right-hand side, as is done when extracting a metric from quantum fidelity. Probability theory Let and be random variables. Then the covariance inequality is given by: After defining an inner product on the set of random variables using the expectation of their product, the Cauchy–Schwarz inequality becomes To prove the covariance inequality using the Cauchy–Schwarz inequality, let and then where denotes variance and denotes covariance. Proofs There are many different proofs of the Cauchy–Schwarz inequality other than those given below. When consulting other sources, there are often two sources of confusion. First, some authors define to be linear in the second argument rather than the first. Second, some proofs are only valid when the field is and not This section gives two proofs of the following theorem: In both of the proofs given below, the proof in the trivial case where at least one of the vectors is zero (or equivalently, in the case where ) is the same. It is presented immediately below only once to reduce repetition. It also includes the easy part of the proof of the Equality Characterization given above; that is, it proves that if and are linearly dependent then By definition, and are linearly dependent if and only if one is a scalar multiple of the other. If where is some scalar then which shows that equality holds in the . The case where for some scalar follows from the previous case: In particular, if at least one of and is the zero vector then and are necessarily linearly dependent (for example, if then where ), so the above computation shows that the Cauchy–Schwarz inequality holds in this case. Consequently, the Cauchy–Schwarz inequality only needs to be proven only for non-zero vectors and also only the non-trivial direction of the Equality Characterization must be shown. Proof via the Pythagorean theorem The special case of was proven above so it is henceforth assumed that Let It follows from the linearity of the inner product in its first argument that: Therefore, is a vector orthogonal to the vector (Indeed, is the projection of onto the plane orthogonal to ) We can thus apply the Pythagorean theorem to which gives The Cauchy–Schwarz inequality follows by multiplying by and then taking the square root. Moreover, if the relation in the above expression is actually an equality, then and hence the definition of then establishes a relation of linear dependence between and The converse was proved at the beginning of this section, so the proof is complete. Proof by analyzing a quadratic Consider an arbitrary pair of vectors . Define the function defined by , where is a complex number satisfying and . Such an exists since if then can be taken to be 1. Since the inner product is positive-definite, only takes non-negative real values. On the other hand, can be expanded using the bilinearity of the inner product: Thus, is a polynomial of degree (unless which is a case that was checked earlier). Since the sign of does not change, the discriminant of this polynomial must be non-positive: The conclusion follows. For the equality case, notice that happens if and only if If then and hence Generalizations Various generalizations of the Cauchy–Schwarz inequality exist. Hölder's inequality generalizes it to norms. More generally, it can be interpreted as a special case of the definition of the norm of a linear operator on a Banach space (Namely, when the space is a Hilbert space). Further generalizations are in the context of operator theory, e.g. for operator-convex functions and operator algebras, where the domain and/or range are replaced by a C-algebra or W-algebra. An inner product can be used to define a positive linear functional. For example, given a Hilbert space being a finite measure, the standard inner product gives rise to a positive functional by Conversely, every positive linear functional on can be used to define an inner product where is the pointwise complex conjugate of In this language, the Cauchy–Schwarz inequality becomes which extends verbatim to positive functionals on C*-algebras: The next two theorems are further examples in operator algebra. This extends the fact when is a linear functional. The case when is self-adjoint, that is, is sometimes known as Kadison's inequality. Another generalization is a refinement obtained by interpolating between both sides of the Cauchy–Schwarz inequality: This theorem can be deduced from Hölder's inequality. There are also non-commutative versions for operators and tensor products of matrices. Several matrix versions of the Cauchy–Schwarz inequality and Kantorovich inequality are applied to linear regression models.	Mathematics	Other algebra topics	null
38173	https://en.wikipedia.org/wiki/Hot%20air%20balloon	Hot air balloon	A hot air balloon is a lighter-than-air aircraft consisting of a bag, called an envelope, which contains heated air. Suspended beneath is a gondola or wicker basket (in some long-distance or high-altitude balloons, a capsule), which carries passengers and a source of heat, in most cases an open flame caused by burning liquid propane. The heated air inside the envelope makes it buoyant, since it has a lower density than the colder air outside the envelope. As with all aircraft, hot air balloons cannot fly beyond the atmosphere. The envelope does not have to be sealed at the bottom, since the air inside the envelope is at about the same pressure as the surrounding air. In modern sport balloons the envelope is generally made from nylon fabric, and the inlet of the balloon (closest to the burner flame) is made from a fire-resistant material such as Nomex. Modern balloons have been made in many shapes, such as rocket ships and the shapes of various commercial products, though the traditional shape is used for most non-commercial and many commercial applications. The hot air balloon is the first successful human-carrying flight technology. The first untethered manned hot air balloon flight in the world was performed in Paris, France, by Jean-François Pilâtre de Rozier and François Laurent d'Arlandes on November 21, 1783, in a balloon created by the Montgolfier brothers. Hot air balloons that can be propelled through the air rather than simply drifting with the wind are known as thermal airships. History Premodern and unmanned balloons A precursor of the hot air balloon was the sky lantern (). Zhuge Liang of the Shu Han kingdom, during the Three Kingdoms era (220–280 CE), used these airborne lanterns for military signaling. The Mongolian army studied Kongming lanterns from China and used them in the Battle of Legnica during the Mongol invasion of Poland in the 13th century. This is the first time ballooning was known in the western world. In the 18th century the Portuguese Jesuit priest Bartolomeu de Gusmão in colonial Brazil envisioned an aerial apparatus named , which was the predecessor of the hot air balloon. The was intended to serve as air vessel in order to facilitate communication and as a strategical device. In 1709 John V of Portugal decided to fund Bartolomeu de Gusmão's project following a petition made by the Jesuit priest, and an unmanned demonstration was performed at Casa da Índia in the presence of John V and the queen, Maria Anna of Austria, with the Italian cardinal Michelangelo Conti, two members of the Portuguese Royal Academy of History, one Portuguese diplomat and one chronicler serving as witnesses. This event would bring some European attention to this event and this project. A later article dated on October 20, 1786, by the London Daily Universal Register would state that the inventor was able to raise himself by the use of his prototype. Also in 1709, the Portuguese Jesuit wrote (Short Manifesto for those who are unaware that is possible to sail through the element air); he also left designs for a manned air vessel. In the 1970s, balloonist Julian Nott hypothesized that the Nazca Lines geoglyphs' creation two millennia ago could have been guided by Nazca leaders in a balloon, possibly the earliest hot air balloon flights in human history. To support this theory, in 1975 he designed and piloted the Nazca Prehistoric Balloon, claiming to have used only methods and materials available to the Pre-Inca Peruvians 1,000 years ago. First manned flight The French brothers Joseph-Michel and Jacques-Étienne Montgolfier developed a hot-air balloon in Annonay, Ardèche, France, and demonstrated it publicly on September 19, 1783, making an unmanned flight lasting 10 minutes. After experimenting with unmanned balloons and flights with animals, the first balloon flight with humans aboard, a tethered flight, performed on or around October 15, 1783, by Jean-Francois Pilatre de Rozier, who made at least one tethered flight from the yard of the Reveillon workshop in the Faubourg Saint-Antoine. Later that same day, Pilatre de Rozier became the second human to ascend into the air, reaching an altitude of , the length of the tether. The first free flight with human passengers was made a few weeks later, on November 21, 1783. King Louis XVI had originally decreed that condemned criminals would be the first pilots, but de Rozier, along with Marquis François d'Arlandes, petitioned successfully for the honor. The first military use of a hot air balloon happened in 1794 during the battle of Fleurus, when the French used the balloon for observation. Modern balloons Modern hot air balloons, with an onboard heat source, were developed by Ed Yost and Jim Winker, beginning during the 1950s; their work resulted in his a first successful flight on October 22, 1960. The first modern hot air balloon to be made in the United Kingdom (UK) was the Bristol Belle, built in 1967. Presently, hot air balloons are used primarily for recreation. Records Hot air balloons are able to fly to extremely high altitudes. On November 26, 2005 Vijaypat Singhania set the world altitude record for highest hot air balloon flight, reaching . He took off from downtown Mumbai, India, and landed south in Panchale. The previous record of had been set by Per Lindstrand on June 6, 1988, in Plano, Texas. On January 15, 1991, the Virgin Pacific Flyer balloon completed the longest flight in a hot air balloon, when Per Lindstrand (born in Sweden, but resident in the UK) and Richard Branson of the UK flew from Japan to Northern Canada. With a volume of 74,000 cubic meters (2.6 million cubic feet), the balloon envelope was the largest ever built for a hot air craft. Designed to fly in the trans-oceanic jet streams, the Pacific Flyer recorded the fastest ground speed for a manned balloon at . The longest duration record was set by Swiss psychiatrist Bertrand Piccard (Auguste Piccard's grandson) and Briton Brian Jones, flying in the Breitling Orbiter 3. It was the first nonstop trip around the world by balloon. The balloon left Château-d'Oex, Switzerland, on March 1, 1999, and landed at 1:02 a.m. on March 21 in the Egyptian desert south of Cairo. The two men exceeded distance, endurance, and time records, traveling 19 days, 21 hours, and 55 minutes. Steve Fossett, flying solo, exceeded the record for briefest time traveling around the world on 3 July 2002 on his sixth attempt, in 320 h 33 min. Fedor Konyukhov flew solo round the world on his first attempt in a hybrid hot air/helium balloon from 11 to 23 July 2016 for a round-the world time of 268 h 20 min. Construction A hot air balloon for manned flight uses a single-layered, fabric gas bag (lifting "envelope"), with an opening at the bottom called the mouth or throat. Attached to the envelope is a basket, or gondola, for carrying the passengers. Mounted above the basket and centered in the mouth is the "burner", which injects a flame into the envelope, heating the air within. The heater or burner is fueled by propane, a liquefied gas stored in pressure vessels, similar to high-pressure forklift cylinders. Envelope Modern hot air balloons are usually made of materials such as ripstop nylon or dacron (a polyester). During the manufacturing process, the material is cut into panels and sewn together, along with structural load tapes that carry the weight of the gondola or basket. The individual sections, which extend from the throat to the crown (top) of the envelope, are known as gores or gore sections. Envelopes can have as few as 4 gores or as many as 24 or more. Envelopes often have a crown ring at their very top. This is a hoop of smooth metal, usually aluminium, and approximately in diameter. Vertical load tapes from the envelope are attached to the crown ring. At the bottom of the envelope the vertical load tapes are sewn into loops that are connected to cables (one cable per load tape). These cables, often referred to as flying wires, are connected to the basket by carabiners. Seams The most common technique for sewing panels together is called the French felled, French fell, or double lap seam. The two pieces of fabric are folded over on each other at their common edge, possibly with a load tape as well, and sewn together with two rows of parallel stitching. Other methods include a flat lap seam, in which the two pieces of fabric are held together simply with two rows of parallel stitching, and a zigzag, where parallel zigzag stitching holds a double lap of fabric. Coatings The fabric (or at least part of it, the top 1/3, for example) may be coated with a sealer, such as silicone or polyurethane, to make it impermeable to air. It is often the degradation of this coating and the corresponding loss of impermeability that ends the effective life of an envelope, not weakening of the fabric itself. Heat, moisture, and mechanical wear-and-tear during set-up and pack-up are the primary causes of degradation. Once an envelope becomes too porous to fly, it may be retired and discarded or perhaps used as a "rag bag": cold-inflated and opened for children to run through. Products for recoating the fabric are becoming available commercially. Sizes and capacity A range of envelope sizes is available. The smallest, one-person, basket-less balloons (called "Hoppers" or "Cloudhoppers") have as little as of envelope volume; for a perfect sphere the radius would be around . At the other end of the scale, balloons used by commercial sightseeing operations may be able to carry well over two dozen people, with envelope volumes of up to . The most-used size is about , allowing to carry 3 to 5 people. Vents The top of the balloon usually has a vent of some sort, enabling the pilot to release hot air to slow an ascent, start a descent, or increase the rate of descent, usually for landing. Some hot air balloons have turning vents, which are side vents that, when opened, cause the balloon to rotate. Such vents are particularly useful for balloons with rectangular baskets, to facilitate aligning the wider side of the basket for landing. The most common type of top vent is a disk-shaped flap of fabric called a parachute vent, invented by Tracy Barnes. The fabric is connected around its edge to a set of "vent lines" that converge in the center. (The arrangement of fabric and lines roughly resembles a parachute—thus the name.) These "vent lines" are themselves connected to a control line that runs to the basket. A parachute vent is opened by pulling on the control line. Once the control line is released, the pressure of the remaining hot air pushes the vent fabric back into place. A parachute vent can be opened briefly while in flight to initiate a rapid descent. (Slower descents are initiated by allowing the air in the balloon to cool naturally.) The vent is pulled open completely to collapse the balloon after landing. An older, and presently less commonly used, style of top vent is called a "Velcro-style" vent. This too is a disk of fabric at the top of the balloon. However, rather than having a set of "vent lines" that can repeatedly open and close the vent, the vent is secured by "hook and loop" fasteners (such as Velcro) and is only opened at the end of the flight. Balloons equipped with a Velcro-style vent typically have a second "maneuvering vent" built into the side (as opposed to the top) of the balloon. Another common type of top design is the "smart vent", which, rather than lowering a fabric disc into the envelope as in the "parachute" type, gathers the fabric together in the center of the opening. This system can theoretically be used for in-flight maneuvering, but is more commonly used only as a rapid-deflation device for use after landing, of particular value in high winds. Other designs, such as the "pop top" and "MultiVent" systems, have also attempted to address the need for rapid deflation on landing, but the parachute top remains popular as an all-around maneuvering and deflation system. Shape Besides special shapes, possibly for marketing purposes, there are several variations on the traditional "inverted tear drop" shape. The simplest, often used by home builders, is a hemisphere on top of a truncated cone. More sophisticated designs attempt to minimize the circumferential stress on the fabric, with different degrees of success depending on whether they take fabric weight and varying air density into account. This shape may be referred to as "natural". Finally, some specialized balloons are designed to minimize aerodynamic drag (in the vertical direction) to improve flight performance in competitions. Basket Hot air balloon baskets are commonly made of woven wicker or rattan. These materials have proven to be sufficiently light, strong, and durable for balloon flight. Such baskets are usually rectangular or triangular in shape. They vary in size from just big enough for two people to large enough to carry thirty. Larger baskets often have internal partitions for structural bracing and to compartmentalize the passengers. Small holes may be woven into the side of the basket to act as foot holds for passengers climbing in or out. Baskets may also be made of aluminium, especially a collapsible aluminium frame with a fabric skin, to reduce weight or increase portability. These may be used by pilots without a ground crew or who are attempting to set altitude, duration, or distance records. Other specialty baskets include the fully enclosed gondolas used for around-the-world attempts and baskets that consist of little more than a seat for the pilot and perhaps one passenger. Burner The burner unit gasifies liquid propane, mixes it with air, ignites the mixture, and directs the flame and exhaust into the mouth of the envelope. Burners vary in power output; each will generally produce 2 to 3 MW of heat (7 to 10 million BTUs per hour), with double, triple, or quadruple burner configurations installed where more power is needed. The pilot actuates a burner by opening a propane valve, known as a blast valve. The valve may be spring-loaded, so that it closes automatically, or it may stay open until closed by the pilot. The burner has a pilot light to ignite the propane and air mixture. The pilot light may be lit by the pilot with an external device, such as a flint striker or a lighter, or with a built-in piezoelectric spark. Where more than one burner is present, the pilot can use one or more at a time, depending on the desired heat output. Each burner has a metal coil of propane tubing the flame shoots through to preheat the incoming liquid propane. The burner unit may be suspended from the mouth of the envelope or supported rigidly over the basket. The burner unit may be mounted on a gimbal to enable the pilot to aim the flame and avoid overheating the envelope fabric. A burner may have a secondary propane valve that releases propane more slowly and thereby generates a different sound. This is called a whisper burner and is used for flight over livestock to lessen the chance of spooking them. It also generates a more yellow flame and is used for night glows because it lights up the inside of the envelope better than the primary valve. Fuel tanks Propane fuel tanks are usually cylindrical pressure vessels made from aluminium, stainless steel, or titanium with a valve at one end to feed the burner and to refuel. They may have a fuel gauge and a pressure gauge. Common tank sizes are . They may be intended for upright or horizontal use and may be mounted inside or outside the basket. The pressure necessary to force the fuel through the line to the burner may be supplied by the vapor pressure of the propane itself, if warm enough, or by the introduction of an inert gas such as nitrogen. Tanks may be preheated with electrical heat tapes to produce sufficient vapor pressure for cold-weather flying. Warmed tanks are usually also wrapped in an insulating blanket to preserve heat during the setup and flight. Instrumentation A balloon may be outfitted with a variety of instruments to aid the pilot. These commonly include an altimeter, a rate-of-climb (vertical-speed) indicator known as a variometer, envelope (air) temperature, and ambient (air) temperature. A GPS receiver can be useful to indicate ground speed (traditional aircraft air-speed indicators would be useless) and direction. Combined mass The combined mass of an average system can be calculated as follows: {\| class="wikitable" \|- ! Component ! Pounds ! Kilograms ! Mass fraction \|- \| envelope \| align="right" \|250 \| align="right" \|113.4 \| \|- \| 5-passenger basket \| align="right" \|140 \| align="right" \|63.5 \| \|- \| Double burner \| align="right" \|50 \| align="right" \|22.7 \| \|- \| 3 fuel tanks full of propane \| align="right" \|3 × 135 = 405 \| align="right" \|183.7 \| \|- \| 5 passengers \| align="right" \|5 × 150 = 750 \| align="right" \|340.2 \| \|- \| Subtotal \| align="right" \|1595 \| align="right" \|723.5 \| \|- \| of heated air* \| align="right" \|5922 \| align="right" \|2686.2 \| \|- \| Total \| align="right" \|(3.76 tons) 7517 \| align="right" \|3409.7 \| \|} * Using a density of for dry air heated to . Theory of operation Generating lift Increasing the air temperature inside the envelope makes it less dense than the surrounding (ambient) air. The balloon floats because of the buoyant force exerted on it. This force is the same force that acts on objects when they are in water and is described by Archimedes' principle. The amount of lift (or buoyancy) provided by a hot air balloon depends primarily upon the difference between the temperature of the air inside the envelope and the temperature of the air outside the envelope. For most envelopes made of nylon fabric, the maximal internal temperature is limited to approximately 120 °C (250 °F). The melting point of nylon is significantly greater than this maximal operating temperature—about 230 °C (450 °F)—but higher temperatures cause the strength of the nylon fabric to degrade more quickly over time. With a maximal operating temperature of 120 °C (250 °F), balloon envelopes can generally be flown for between 400 and 500 hours before the fabric needs to be replaced. Many balloon pilots operate their envelopes at temperatures significantly less than the maximum to extend envelope-fabric life. The lift generated by of dry air heated to various temperatures may be calculated as follows: {\| class="wikitable" \|- ! Air temperature ! Air density ! Air mass ! Lift generated \|- \| \| \| \| 0 lb, 0 kg \|- \| \| \| \| \|- \| \| \| \| \|} The density of air at is about . The total lift for a balloon of heated to would be . This is just enough to generate neutral buoyancy for the total system mass (not including the heated air trapped in the envelope, of course) stated in the previous section. Liftoff would require a slightly greater temperature, depending on the desired rate of climb. In reality, the air contained in the envelope is not all at the same temperature, as the accompanying thermal image shows, and so these calculations are based on averages. For typical atmospheric conditions (), a hot air balloon heated to requires about 3.91 m3 of envelope volume to lift 1 kilogram (equivalently, 62.5 cu ft/lb). The precise amount of lift provided depends not only upon the internal temperature mentioned above, but the external temperature, altitude above sea level, and humidity of the air surrounding. On a warm day, a balloon cannot lift as much as on a cool day, because the temperature required for launch will exceed the maximum sustainable for nylon envelope fabric. Also, in the lower atmosphere, the lift provided by a hot air balloon decreases about 3% per 1,000 m (1% per 1,000 ft) of altitude gained. Types of Hot Air Balloons There are several different types of hot air balloons, all with different means of taking and sustaining flight. Montgolfier Standard hot air balloons are known as Montgolfier balloons and rely solely on the buoyancy of hot air provided by the burner and contained by the envelope. This style of balloon was developed by the Montgolfier brothers and had its first public demonstration on 4 June 1783 with an unmanned flight lasting 10 minutes, followed later that year with manned flights. Gas Instead of using regular air it is also possible to use lighter than air gasses such as Helium or Hydrogen to lift the balloon, though this means it is technically not a hot air balloon, though they did influence the design of hybrid balloons. Hybrid The 1785 Rozière balloon, is the main type of hybrid balloon, named after its creator, Jean-François Pilâtre de Rozier. It has a separate cell for a lighter-than-air gas (typically helium), as well as a cone below for hot air (as is used in a hot air balloon) to heat the helium at night. Hydrogen gas was used in the very early stages of development but was quickly abandoned due to the danger of introducing an open flame near the gas, for example when Rozier attempted to cross the English Channel with his prototype, the fire used to heat the air ignited the Hydrogen and killed both him and his copilot thirty minutes after takeoff. As such, all modern hybrid balloons now use helium as a lifting gas. These balloons are commonly used for high performance records for hot air balloons. Solar Solar balloons are hot air balloons that use just solar energy captured by an envelope. These envelopes are more specialized than for other hot air balloons, trying to maximize the amount of solar energy they collect. This can be accomplished by rotating the envelope during flight or by having the envelope colored black or another dark color. They were pioneered in the 1970s in Europe by Tracy Barnes, Dominic Michaelis, and in the US by Frederick Espoo and Paul Woessher. Thermal Airship A Thermal airship, or blimp, became a reality in the 1960s. Thermal airships were the first steerable air buoyant vehicles. They utilized tail fins and a rudder and contained strictly hot air rather than a mix with hydrogen or helium. Observation Balloon Observation balloons were deployed as early as the American Civil War and used as reconnaissance towers. The first military funded balloon in America was designed by Thaddeus Lowe on August 2, 1861, for the Union. His design utilized gas from municipal lines to inflate the balloon as he did not have access to a portable generator. Observation balloons during this time were all made using multicolored-silk, wicker baskets, and were vertically oriented and tear shaped. Hydrogen, or illumination gas became the most used inflation fuel by the 20th Century, as it was lighter than air. Observation balloon usage skyrocketed in Britain by the Royal Engineers at the end of the 19th Century, deploying to Sudan in 1885 and to South Africa during the Second Boer War from 1899 to 1902. Steering Due to the overall design of hot air balloons, controlled and precise steering of hot air balloons is not possible; it is possible for pilots to try to achieve basic directional control by changing altitude and catching different wind streams. Wind in the northern hemisphere tends to turn east due to coriolis effect as the altitude increases. Landing The most effective way of landing a hot air balloon is to reduce the energy in the envelope, either by turning down the flame in Montgolfier and Hybrid balloons, or more directly by opening a flap in the envelope that will release the air/gas inside. Safety equipment To help ensure the safety of pilot and passengers, a hot air balloon may carry several pieces of safety equipment. Basket To relight the burner if the pilot light goes out and the optional piezo ignition fails, the pilot should have ready access to a means of backup ignition, such as a flint spark lighter. Many systems, especially those that carry passengers, have completely duplicate fuel and burner systems: two fuel tanks, connected to two separate hoses, which feed two distinct burners. This enables a safe landing in the case of a blockage somewhere in one system or if a system must be disabled because of a fuel leak. A fire extinguisher suitable for extinguishing propane fires is useful. Most balloons carry a 1 or 2 kg AB:E type fire extinguisher. A handling or drop line is mandatory safety equipment in many countries. This is a rope or webbing of 20–30 meters in length attached to the balloon basket with a quick-release connection at one end. In very calm winds the balloon pilot can throw the handling line from the balloon so that the ground crew can guide the balloon safely away from obstructions on the ground. For commercial passenger balloons, a pilot restraint harness is mandatory in some countries. This consists of a hip belt and a webbing line that together allow for some movement while preventing the pilot from being ejected from the basket during a hard landing. Further safety equipment may include a first-aid kit, a fire blanket and a strong rescue knife. Occupants At a minimum, the pilot should wear leather or flame-retardant fiber (such as nomex) gloves, so that they may shut off a gas valve in the case of a leak, even if there is a flame present; quick action in this regard can turn a potential catastrophe into a mere inconvenience. The pilot should additionally wear flame-resistant clothing covering their arms and legs; either natural fiber, such as cotton, linen, hemp, or wool, or engineered flame-retardant fiber, such as nomex, is acceptable in this capacity. Most engineered fibers (with the exception of rayon, which is also safe to wear) are thermoplastic; many are also hydrocarbons. This makes such fabrics very much unsuitable to wear near high temperatures, since non-flame-retardant thermoplastics will melt onto the wearer, and most hydrocarbons, whether fibrous or not, are suitable to use as fuels. Natural fiber will singe rather than melt or burn readily, and flame-retardant fiber generally has a very high melting point and is intrinsically non-flammable. Many pilots also advise their passengers to wear similar protective clothing that covers their arms and legs, as well as strong shoes or boots that offer good ankle support. Finally, some balloon systems, especially those that hang the burner from the envelope instead of supporting it rigidly from the basket, require the use of helmets by the pilot and passengers. Ground crew The ground crew should wear gloves whenever there is a possibility of handling ropes or lines. The mass and exposed surface to air movement of a medium-sized balloon is sufficient to cause rope friction burns to the hands of anyone trying to stop or prevent movement. The ground crew should also wear sturdy shoes and at least long pants in case of the need to access a landing or landed balloon in rough or overgrown terrain. Maintenance and repair As with aircraft, hot air balloons require regular maintenance to remain airworthy. As aircraft made of fabric and that lack direct horizontal control, hot air balloons may occasionally require repairs to rips or snags. While some operations, such as cleaning and drying, may be performed by the owner or pilot, other operations, such as sewing, must be performed by a qualified repair technician and recorded in the balloon's maintenance log book. Maintenance To ensure long life and safe operation, the envelope should be kept clean and dry. This prevents mold and mildew from forming on the fabric and abrasion from occurring during packing, transport, and unpacking due to contact with foreign particles. In the event of a landing in a wet (because of precipitation or early morning or late evening dew) or muddy location (farmer's field), the envelope should be cleaned and laid out or hung to dry. The burner and fuel system must also be kept clean to ensure safe operation on demand. Damaged fuel hoses need to be replaced. Stuck or leaky valves must be repaired or replaced. The wicker basket may require occasional refinishing or repair. The skids on its bottom may require occasional replacement. Balloons in most parts of the world are maintained in accordance with a fixed manufacturer's maintenance schedule that includes regular (100 flight hours or 12 month) inspections, in addition to maintenance work to correct any damage. In Australia, balloons used for carrying commercial passengers must be inspected and maintained by approved workshops. Repair In the case of a snag, burn, or rip in the envelope fabric, a patch may be applied or the affected panel completely replaced. Patches may be held in place with glue, tape, stitching, or a combination of these techniques. Replacing an entire panel requires the stitching around the old panel to be removed, and a new panel to be sewn in with the appropriate technique, thread, and stitch pattern. Licensing Depending on the size of the balloon, location, and intended use, hot air balloons and their pilots need to comply with a variety of regulations. Balloons As with other aircraft in the US, balloons must be registered (have an N-number), have an airworthiness certificate, and pass annual inspections. Balloons below a certain size (empty weight of less than 155 pounds or 70 kg including envelope, basket, burners and empty fuel tanks) can be used as an ultralight aircraft. Pilots In Australia In Australia, private balloon pilots are managed by the Australian Ballooning Federation and typically become members of regional hot air ballooning clubs. Commercial operations carrying fare paying passengers or charging for promotional flights must operate under an Air Operators Certificate from the Australian Civil Aviation and Safety Authority (CASA) with a nominated Chief Pilot. Pilots must have different degrees of experience before they are allowed to progress to larger balloons. Hot air balloons must be registered aircraft with CASA and are subject to regular airworthiness checks by authorised personnel. In the UK In the UK, the person in command must hold a valid Private Pilot's Licence issued by the Civil Aviation Authority specifically for ballooning; this is known as the PPL(B). There are two types of commercial balloon licences: CPL(B) Restricted and CPL(B) (Full). The CPL(B) Restricted is required if the pilot is undertaking work for a sponsor or being paid by an external agent to operate a balloon. The pilot can fly a sponsored balloon with everything paid for with a PPL unless asked to attend any event. Then a CPL(B) Restricted is required. The CPL(B) is required if the pilot is flying passengers for money. The balloon then needs a transport category C of A (certificate of air worthiness). If the pilot is only flying sponsor's guests and not charging money for flying other passengers, then the pilot is exempted from holding an AOC (air operator's certificate) though a copy of it is required. For passenger flying the balloon also requires a maintenance log. In the United States In the United States, a pilot of a hot air balloon must have a pilot certificate from the Federal Aviation Administration (FAA), carrying the rating of "Lighter-than-air free balloon", and unless the pilot is also qualified to fly gas balloons, will also carry this limitation: "Limited to hot-air balloons with airborne heater". A pilot does not need a license to fly an ultralight aircraft, but training is highly advised, and some hot-air balloons meet the criteria. To carry paying passengers for hire (and attend some balloon festivals), a pilot must have a commercial pilot certificate. Commercial hot air balloon pilots may also act as hot air balloon flight instructors. While most balloon pilots fly for the pure joy of floating through the air, many are able to make a living as a professional balloon pilot. Some professional pilots fly commercial passenger sightseeing flights, while others fly corporate advertising balloons. Accidents and incidents 1989 Alice Springs hot air balloon crash: On 13 August 1989, two hot air balloons collided at Alice Springs, Northern Territory, Australia, causing one to fall, killing all 13 people on board. 2011 Somerset hot air balloon crash: On 1 January 2011, a hot air balloon attempting a high-altitude flight crashed at Pratten's Bowls Club in Westfield, Somerset, near Bath, England, killing both people on board. 2012 Carterton hot air balloon crash: On 7 January 2012, a hot air balloon collided with a power line, caught fire and crashed at Carterton, North Island, New Zealand, killing all 11 people on board. 2012 Ljubljana Marshes hot air balloon crash: On 23 August 2012, a storm blew a hot air balloon to the ground, causing it to catch fire on impact near Ljubljana, Slovenia. The crash killed 6 of the 32 people on board, and injured the other 26. 2013 Luxor hot air balloon crash: On 26 February 2013, a hot air balloon carrying foreign tourists ignited and crashed near the ancient city of Luxor, Egypt, killing 19 of the 21 people on board, making it the deadliest balloon accident in history. 2016 Lockhart hot air balloon crash: On 30 July 2016, a hot air balloon carrying 16 people caught fire and crashed near Lockhart, Texas. There were no survivors. 2021 Albuquerque hot air balloon crash: On 26 June 2021, a hot air balloon carrying five people made contact with a power line and crashed in Albuquerque, New Mexico. All five people on board died as a result of the accident. On 14 January 2024, a hot air balloon crashed outside Eloy, Arizona, killing the three passengers and pilot. Eight skydivers had exited the balloon immediately prior to the incident. Manufacturers The largest manufacturer of hot air balloons is Cameron Balloons of Bristol, England, which also owns Lindstrand Balloons of Oswestry, England. Cameron Balloons, Lindstrand Balloons and another English balloon manufacturing company Thunder and Colt (since acquired by Cameron) have been innovators and developers of special shaped balloons. These hot-air balloons use the same principle of lift as conventional inverted teardrop-shaped balloons, but often sections of the special balloon envelope shape do not contribute to the balloon's ability to stay aloft. The second largest manufacturer of hot air balloons is Ultramagic company, based in Spain, which produces from 80 to 120 balloons per year. Ultramagic can produce very large balloons, such as the N-500 that accommodates as many as 27 persons in the basket, and has also produced many balloons with special shapes, as well as cold-air inflatables. One of the three largest companies in the world is Kubicek Balloons. From its factory in Brno, Czechia the company ships its products worldwide. Produces from 100 to 115 balloons per year. Kubicek company also focus on special shape balloons, FAA/EASA type certified and are delivered with a Standard Airworthiness Certificate. In the USA Aerostar International, Inc. of Sioux Falls, South Dakota was North America's largest balloon manufacturer and a close second in world manufacturing before ceasing to build balloons in January 2007. The oldest U.S. certified manufacturer is now Adams Balloons out of Albuquerque, New Mexico. Firefly Balloons, formerly The Balloon Works, is a manufacturer of hot air balloons in Statesville, North Carolina. Another manufacturer is Head Balloons, Inc. of Helen, Georgia. The major manufacturers in Canada are Sundance Balloons and Fantasy Sky Promotions. Other manufacturers include Kavanagh Balloons of Australia, Schroeder Fire Balloons of Germany, Kubicek Balloons of the Czech Republic, and LLopis Balloons of France.	Technology	Aviation	null
38198	https://en.wikipedia.org/wiki/Pioneer%2010	Pioneer 10	Pioneer 10 (originally designated Pioneer F) is a NASA space probe launched in 1972 that completed the first mission to the planet Jupiter. Pioneer 10 became the first of five planetary probes and 11 artificial objects to achieve the escape velocity needed to leave the Solar System. This space exploration project was conducted by the NASA Ames Research Center in California. The space probe was manufactured by TRW Inc. Pioneer 10 was assembled around a hexagonal bus with a diameter parabolic dish high-gain antenna, and the spacecraft was spin stabilized around the axis of the antenna. Its electric power was supplied by four radioisotope thermoelectric generators that provided a combined 155 watts at launch. It was launched on March 3, 1972, at 01:49:00 UTC (March 2 local time), by an Atlas-Centaur rocket from Cape Canaveral, Florida. Between July 15, 1972, and February 15, 1973, it became the first spacecraft to traverse the asteroid belt. Photography of Jupiter began on November 6, 1973, at a range of , and about 500 images were transmitted. The closest approach to the planet was on December 3, 1973, at a range of . During the mission, the on-board instruments were used to study the asteroid belt, the environment around Jupiter, the solar wind, cosmic rays, and eventually the far reaches of the Solar System and heliosphere. Radio communications were lost with Pioneer 10 on January 23, 2003, because of the loss of electric power for its radio transmitter. At the time, the probe was from Earth. Mission background History In the 1960s, American aerospace engineer Gary Flandro of the NASA Jet Propulsion Laboratory conceived of a mission, known as the Planetary Grand Tour, that would exploit a rare alignment of the outer planets of the Solar System. This mission would ultimately be accomplished in the late 1970s by the two Voyager probes, but in order to prepare for it, NASA decided in 1964 to experiment with launching a pair of probes to the outer Solar System. An advocacy group named the Outer Space Panel and chaired by American space scientist James A. Van Allen, worked out the scientific rationale for exploring the outer planets. NASA Goddard Spaceflight Center put together a proposal for a pair of "Galactic Jupiter Probes" that would pass through the asteroid belt and visit Jupiter. These were to be launched in 1972 and 1973 during favorable windows that occurred only a few weeks every 13 months. Launch during other time intervals would have been more costly in terms of propellant requirements. Approved by NASA in February 1969, the twin spacecraft were designated Pioneer F and Pioneer G before launch; later, they were named Pioneer 10 and Pioneer 11 respectively. They formed part of the Pioneer program, a series of United States uncrewed space missions launched between 1958 and 1978. This model was the first in the series to be designed for exploring the outer Solar System. Based on proposals issued throughout the 1960s, the early mission objectives were to explore the interplanetary medium past the orbit of Mars, study the asteroid belt and assess the possible hazard to spacecraft traveling through the belt, and explore Jupiter and its environment. Later development-stage objectives included the probe closely approaching Jupiter to provide data on the effect the environmental radiation surrounding Jupiter would have on the spacecraft instruments. More than 150 scientific experiments were proposed for the missions. The experiments to be carried on the spacecraft were selected in a series of planning sessions during the 1960s, then were finalized by early 1970. These would be to perform imaging and polarimetry of Jupiter and several of its satellites, make infrared and ultraviolet observations of Jupiter, detect asteroids and meteoroids, determine the composition of charged particles, and to measure magnetic fields, plasma, cosmic rays and the zodiacal light. Observation of the spacecraft communications as it passed behind Jupiter would allow measurements of the planetary atmosphere, while tracking data would improve estimates of the mass of Jupiter and its moons. NASA Ames Research Center, rather than Goddard, was selected to manage the project as part of the Pioneer program. The Ames Research Center, under the direction of Charles F. Hall, was chosen because of its previous experience with spin-stabilized spacecraft. The requirements called for a small, lightweight spacecraft which was magnetically clean and which could perform an interplanetary mission. It was to use spacecraft modules that had already been proven in the Pioneer 6 through 9 missions. Ames commissioned a documentary film by George Van Valkenburg titled Jupiter Odyssey. It received numerous international awards, and is visible on Van Valkenburg's YouTube channel. In February 1970, Ames awarded a combined US$380 million contract to TRW Inc. for building both of the Pioneer 10 and 11 vehicles, bypassing the usual bidding process to save time. B. J. O'Brien and Herb Lassen led the TRW team that assembled the spacecraft. Design and construction of the spacecraft required an estimated 25 million man-hours. An engineer from TRW quipped, "This spacecraft is guaranteed for two years of interplanetary flight. If any component fails within that warranty period, just return the spacecraft to our shop and we will repair it free of charge." To meet the schedule, the first launch would need to take place between February 29 and March 17 so that it could arrive at Jupiter in November 1974. This was later revised to an arrival date of December 1973 in order to avoid conflicts with other missions over the use of the Deep Space Network for communications, and to miss the period when Earth and Jupiter would be at opposite sides of the Sun. The encounter trajectory for Pioneer 10 was selected to maximize the information returned about the radiation environment around Jupiter, even if this caused damage to some systems. It would come within about three times the radius of the planet, which was thought to be the closest it could approach and still survive the radiation. The trajectory chosen would give the spacecraft a good view of the sunlit side. Spacecraft design The Pioneer 10 bus measures deep and with six long panels forming the hexagonal structure. The bus houses propellant to control the orientation of the probe and eight of the eleven scientific instruments. The equipment compartment lay within an aluminum honeycomb structure to provide protection from meteoroids. A layer of insulation, consisting of aluminized mylar and kapton blankets, provides passive thermal control. Heat was generated by the dissipation of 70 to 120 watts (W) from the electrical components inside the compartment. The heat range was maintained within the operating limits of the equipment by means of louvers located below the mounting platform. The spacecraft had a launch mass of about . At launch, the spacecraft carried of liquid hydrazine monopropellant in a diameter spherical tank. Orientation of the spacecraft is maintained with six 4.5 N, hydrazine thrusters mounted in three pairs. Pair one maintained a constant spin-rate of 4.8 rpm, pair two controlled the forward thrust, and pair three controlled the attitude. The attitude pair were used in conical scanning maneuvers to track Earth in its orbit. Orientation information was also provided by a star sensor able to reference Canopus, and two Sun sensors. Power and communications Pioneer 10 uses four SNAP-19 radioisotope thermoelectric generators (RTGs). They are positioned on two three-rod trusses, each in length and 120 degrees apart. This was expected to be a safe distance from the sensitive scientific experiments carried on board. Combined, the RTGs provided 155 W at launch, and decayed to 140 W in transit to Jupiter. The spacecraft required 100 W to power all systems. The generators are powered by the radioisotope fuel plutonium-238, which is housed in a multi-layer capsule protected by a graphite heat shield. The pre-launch requirement for the SNAP-19 was to provide power for two years in space; this was greatly exceeded during the mission. The plutonium-238 has a half-life of 87.74 years, so that after 29 years the radiation being generated by the RTGs was at 80% of its intensity at launch. However, steady deterioration of the thermocouple junctions led to a more rapid decay in electrical power generation, and by 2001 the total power output was 65 W. As a result, later in the mission only selected instruments could be operated at any one time. The space probe includes a redundant system of transceivers, one attached to the narrow-beam, high-gain antenna, the other to an omni-antenna and medium-gain antenna. The parabolic dish for the high-gain antenna is in diameter and made from an aluminum honeycomb sandwich material. The spacecraft was spun about an axis that is parallel to the axis of this antenna so that it could remain oriented toward the Earth. Each transceiver is an 8 W one and transmits data across the S-band using 2110 MHz for the uplink from Earth and 2292 MHz for the downlink to Earth with the Deep Space Network tracking the signal. Data to be transmitted is passed through a convolutional encoder so that most communication errors could be corrected by the receiving equipment on Earth. The data transmission rate at launch was 256 bit/s, with the rate degrading by about 1.27 millibit/s for each day during the mission. Much of the computation for the mission is performed on Earth and transmitted to the spacecraft, where it was able to retain in memory up to five commands of the 222 possible entries by ground controllers. The spacecraft includes two command decoders and a command distribution unit, a very limited form of a processor, to direct operations on the spacecraft. This system requires that mission operators prepare commands long in advance of transmitting them to the probe. A data storage unit is included to record up to 6,144 bytes of information gathered by the instruments. The digital telemetry unit is used to prepare the collected data in one of the thirteen possible formats before transmitting it back to Earth. Scientific instruments Mission profile Launch and trajectory Pioneer 10 was launched on March 3, 1972, at 01:49:00 UTC (8:49 p.m. Eastern Standard Time on March 2) by the National Aeronautics and Space Administration from Space Launch Complex 36A in Florida, aboard an Atlas-Centaur. The third stage of the expendable vehicle consisted of a solid fuel Star-37E stage (TE-M-364-4) developed specifically for the Pioneer missions. This stage provided about of thrust and spun up the spacecraft. The spacecraft had an initial spin rate of 30 rpm. Twenty minutes following the launch, the vehicle's three booms were extended, which slowed the rotation rate to 4.8 rpm. This rate was maintained throughout the voyage. The launch vehicle accelerated the probe for a net interval of 17 minutes, reaching a velocity of . After the high-gain antenna was contacted, several of the instruments were activated for testing while the spacecraft was moving through the Earth's radiation belts. Ninety minutes after launch, the spacecraft reached interplanetary space. Pioneer 10 passed by the Moon in 11 hours, and became the fastest human-made object at that time. Two days after launch, the scientific instruments were turned on, beginning with the cosmic ray telescope. After ten days, all of the instruments were active. During the first seven months of the journey, the spacecraft made three course corrections. The on-board instruments underwent checkouts, with the photometers examining Jupiter and the Zodiacal light, and experiment packages being used to measure cosmic rays, magnetic fields and the solar wind. The only anomaly during this interval was the failure of the Canopus sensor, which instead required the spacecraft to maintain its orientation using the two Sun sensors. While passing through interplanetary medium, Pioneer 10 became the first mission to detect interplanetary atoms of helium. It also observed high-energy ions of aluminum and sodium in the solar wind. The spacecraft recorded important heliophysics data in early August 1972 by registering a solar shock wave when it was at a distance of . On July 15, 1972, Pioneer 10 was the first spacecraft to enter the asteroid belt, located between the orbits of Mars and Jupiter. The project planners expected a safe passage through the belt, and the closest the trajectory would take the spacecraft to any of the known asteroids was . One of the nearest approaches was to the asteroid 307 Nike on December 2, 1972. The on-board experiments demonstrated a deficiency of particles below a micrometer (μm) in the belt, as compared to the vicinity of the Earth. The density of dust particles between 10 and 100 μm did not vary significantly during the trip from the Earth to the outer edge of the belt. Only for particles with a diameter of 100 μm to 1.0 mm did the density show an increase, by a factor of three in the region of the belt. No fragments larger than a millimeter were observed in the belt, indicating these are likely rare; certainly much less common than anticipated. As the spacecraft did not collide with any particles of substantial size, it passed safely through the belt, emerging on the other side about February 15, 1973. Encounter with Jupiter On November 6, 1973, the Pioneer 10 spacecraft was at a distance of from Jupiter. Testing of the imaging system began, and the data were successfully received back at the Deep Space Network. A series of 16,000 commands were then uploaded to the spacecraft to control the flyby operations during the next sixty days. The orbit of the outer moon Sinope was crossed on November 8. The bow shock of Jupiter's magnetosphere was reached on November 16, as indicated by a drop in the velocity of the solar wind from to . The magnetopause was passed through a day later. The spacecraft instruments confirmed that the magnetic field of Jupiter was inverted compared to that of Earth. By the 29th, the orbits of all of the outermost moons had been passed, and the spacecraft was operating flawlessly. Red and blue pictures of Jupiter were being generated by the imaging photopolarimeter as the rotation of the spacecraft carried the instrument's field of view past the planet. These red and blue colors were combined to produce a synthetic green image, allowing a three-color combination to produce the rendered image. On November 26, a total of twelve such images were received back on Earth. By December 2, the image quality exceeded the best images made from Earth. These were being displayed in real-time back on Earth, and the Pioneer program would later receive an Emmy award for this presentation to the media. The motion of the spacecraft produced geometric distortions that later had to be corrected by computer processing. During the encounter, a total of more than 500 images were transmitted. The trajectory of the spacecraft took it along the magnetic equator of Jupiter, where the ion radiation was concentrated. Peak flux for this electron radiation is 10,000 times stronger than the maximum radiation around the Earth. Pioneer 10 passed through the inner radiation belts within 20 RJ, receiving an integrated dose of 200,000 rads from electrons and 56,000 rads from protons (in comparison, a whole body dose of 500 rads is fatal to humans). The level of radiation at Jupiter was ten times more powerful than Pioneer's designers had predicted, leading to fears that the probe would not survive. Starting on December 3, the radiation around Jupiter caused false commands to be generated. Most of these were corrected by contingency commands, but an image of Io and a few close-ups of Jupiter were lost. Similar false commands would be generated on the way out from the planet. Nonetheless, Pioneer 10 did succeed in obtaining images of the moons Ganymede and Europa. The image of Ganymede showed low albedo features in the center and near the south pole, while the north pole appeared brighter. Europa was too far away to obtain a detailed image, although some albedo features were apparent. The trajectory of Pioneer 10 was chosen to take it behind Io, allowing the refractive effect of the moon's atmosphere on the radio transmissions to be measured. This demonstrated that the ionosphere of the moon was about above the surface of the day side, and the density ranged from 60,000 electrons per cubic centimeter on the dayside to 9,000 electrons per cubic centimeter on the night face. An unexpected discovery was that Io was orbiting within a cloud of hydrogen that extended for about , with a width and height of . A smaller, cloud was believed to have been detected near Europa. It was not until after Pioneer 10 had cleared the asteroid belt that NASA selected a trajectory towards Jupiter which included a slingshot effect to send the spacecraft out of the Solar System. Pioneer 10 was the first spacecraft to attempt such a maneuver, a model for future missions. Such an extended mission was not in the initial proposal, but was planned for prior to launch. At the closest approach, the velocity of the spacecraft reached , and it came within of the outer atmosphere of Jupiter. Close-up images of the Great Red Spot and the terminator were obtained. Communication with the spacecraft then ceased as it passed behind the planet. The radio occultation data allowed the temperature structure of the outer atmosphere to be measured, showing a temperature inversion between the altitudes with 10 and 100 mbar pressures. Temperatures at the 10-mbar level ranged from , while temperatures at the 100 mbar level were . The spacecraft generated an infrared map of the planet, which confirmed the idea that the planet radiated more heat than it received from the Sun. Crescent images of the planet were then returned as Pioneer 10 moved away from the planet. As the spacecraft headed outward, it again passed the bow shock of Jupiter's magnetosphere. As this front is constantly shifting in space because of dynamic interaction with the solar wind, the vehicle crossed the bow shock a total of 17 times before it escaped completely. Deep space Pioneer 10 crossed the orbit of Saturn in 1976 and the orbit of Uranus in 1979. On June 13, 1983, the craft crossed the orbit of Neptune, and so became the first human-made object to leave the proximity of the major planets of the Solar System. The mission came to an official end on March 31, 1997, when it had reached a distance of from the Sun, though the spacecraft was still able to transmit coherent data after this date. After March 31, 1997, Pioneer 10s weak signal continued to be tracked by the Deep Space Network to aid the training of flight controllers in the process of acquiring deep-space radio signals. There was an Advanced Concepts study applying chaos theory to extract coherent data from the fading signal. The last successful reception of telemetry was received from Pioneer 10 on April 27, 2002; subsequent signals were barely strong enough to detect and provided no usable data. The final, very weak signal from Pioneer 10 was received on January 23, 2003, when it was from Earth. Further attempts to contact the spacecraft were unsuccessful. A final attempt was made on the evening of March 4, 2006, the last time the antenna would be correctly aligned with Earth. No response was received from Pioneer 10. NASA decided that the RTG units had probably fallen below the power threshold needed to operate the transmitter. Hence, no further attempts at contact were made. Timeline Current status and future On July 18, 2023, Voyager 2 overtook Pioneer 10, making Pioneer 10 the third farthest spacecraft from the Sun after Voyager 1 and Voyager 2. As of June 2024, the probe is estimated to be from Earth and from the Sun. Sunlight takes 18.9 hours to reach Pioneer 10. The brightness of the Sun from the spacecraft is magnitude −16.0. Pioneer 10 is currently travelling in the direction of the constellation Taurus. If left undisturbed, Pioneer 10 and its sister craft Pioneer 11 will join the two Voyager spacecraft and the New Horizons spacecraft in leaving the Solar System to wander the interstellar medium. The Pioneer 10 trajectory is expected to take it in the general direction of the star Aldebaran, currently located at a distance of about 68 light years. If Aldebaran had zero relative velocity, it would require more than two million years for the spacecraft to reach it. Well before that, in about 90,000 years, Pioneer 10 will pass about from the late K-type star HIP 117795. This is the closest stellar flyby in the next few million years of all the Pioneer, Voyager, and New Horizons spacecraft, which are leaving the Solar System. A backup unit, Pioneer H, is currently on display in the "Milestones of Flight" gallery at the National Air and Space Museum in Washington, D.C. Many elements of the mission proved to be critical in the planning of the Voyager program. Pioneer plaque Because it was strongly advocated by Carl Sagan, Pioneer 10 and Pioneer 11 carry a gold-anodized aluminum plaque in case either spacecraft is ever found by intelligent life-forms from another planetary system. The plaques feature the nude figures of a human male and female along with several symbols that are designed to provide information about the origin of the spacecraft. The plaque is attached to the antenna support struts where it would be shielded from interstellar dust. Pioneer 10 in popular media In the film Star Trek V: The Final Frontier, a Klingon Bird-of-Prey destroys Pioneer 10 as target practice. In the serialized speculative fiction multimedia narrative 17776, one of the main characters is a sentient Pioneer 10.	Technology	Unmanned spacecraft	null
38203	https://en.wikipedia.org/wiki/Menstruation	Menstruation	Menstruation (also known as a period, among other colloquial terms) is the regular discharge of blood and mucosal tissue from the inner lining of the uterus through the vagina. The menstrual cycle is characterized by the rise and fall of hormones. Menstruation is triggered by falling progesterone levels, and is a sign that pregnancy has not occurred. The first period, a point in time known as menarche, usually begins between the ages of 12 and 15. Menstruation starting as young as 8 years would still be considered normal. The average age of the first period is generally later in the developing world, and earlier in the developed world. The typical length of time between the first day of one period and the first day of the next is 21 to 45 days in young women. In adults, the range is between 21 and 35 days with the average being 28 days. Bleeding usually lasts around 2 to 7 days. Periods stop during pregnancy and typically do not resume during the initial months of breastfeeding. Lochia occurs after childbirth. Menstruation, and with it the possibility of pregnancy, ceases after menopause, which usually occurs between 45 and 55 years of age. Up to 80% of women do not experience problems sufficient to disrupt daily functioning either during menstruation or in the days leading up to menstruation. Symptoms in advance of menstruation that do interfere with normal life are called premenstrual syndrome (PMS). Some 20 to 30% of women experience PMS, with 3 to 8% experiencing severe symptoms. These include acne, tender breasts, bloating, feeling tired, irritability, and mood changes. Other symptoms some women experience include painful periods (estimates are between 50 and 90%) and heavy bleeding during menstruation and abnormal bleeding at any time during the menstrual cycle. A lack of periods, known as amenorrhea, is when periods do not occur by age 15 or have not re-occurred in 90 days. Characteristics Length and duration The first menstrual period occurs after the onset of pubertal growth, and is called menarche. The average age of menarche is 12 to 15 years. However, it may occur as early as eight. The average age of the first period is generally later in the developing world, and earlier in the developed world. The average age of menarche has changed little in the United States since the 1950s. Menstruation is the most visible phase of the menstrual cycle and its beginning is used as the marker between cycles. The first day of menstrual bleeding is the date used for the last menstrual period (LMP). The typical length of time between the first day of one period and the first day of the next is 21 to 45 days in young women, and 21 to 35 days in adults. The average length is 28 days; one study estimated it at 29.3 days. The variability of menstrual cycle lengths is highest for women under 25 years of age and is lowest, that is, most regular, for ages 25 to 39 years. The variability increases slightly for women aged 40 to 44 years. Perimenopause is when a woman's fertility declines, and menstruation occurs less regularly in the years leading up to the final menstrual period, when a woman stops menstruating completely and is no longer fertile. The medical definition of menopause is one year without a period and typically occurs between 45 and 55 years in Western countries. Menopause before age 45 is considered premature in industrialized countries. Illnesses, certain surgeries, or medical treatments may cause menopause to occur earlier than it might have otherwise. Bleeding The average volume of menstrual fluid during a monthly menstrual period is with considered typical. Menstrual fluid is the correct name for the flow, although many people prefer to refer to it as menstrual blood. Menstrual fluid is reddish-brown, a slightly darker color than venous blood. About half of menstrual fluid is blood. This blood contains sodium, calcium, phosphate, iron, and chloride, the extent of which depends on the woman. As well as blood, the fluid consists of cervical mucus, vaginal secretions, and endometrial tissue. Vaginal fluids in menses mainly contribute water, common electrolytes, organ moieties, and at least 14 proteins, including glycoproteins. Many women and girls notice blood clots during menstruation. These appear as clumps of blood that may look like tissue. If there was a miscarriage or a stillbirth, examination under a microscope can confirm if it was endometrial tissue or pregnancy tissue (products of conception) that was shed. Sometimes menstrual clots or shed endometrial tissue is incorrectly thought to indicate an early-term miscarriage of an embryo. An enzyme called plasmin – contained in the endometrium – tends to inhibit the blood from clotting. The amount of iron lost in menstrual fluid is relatively small for most women. In one study, premenopausal women who exhibited symptoms of iron deficiency were given endoscopies. 86% of them actually had gastrointestinal disease and were at risk of being misdiagnosed simply because they were menstruating. Heavy menstrual bleeding, occurring monthly, can result in anemia. Hormonal changes Side effects Menstrual health overview Moods and premenstrual syndrome (PMS) Cramps In most women, various physical changes are brought about by fluctuations in hormone levels during the menstrual cycle. This includes muscle contractions of the uterus (menstrual cramping) that can precede or accompany menstruation. Many women experience painful cramps, also known as dysmenorrhea, during menstruation. Among adult women, that pain is severe enough to affect daily activity in only 2%–28%. Severe symptoms that disrupt daily activities and functioning may be diagnosed as premenstrual dysphoric disorder. These symptoms can be severe enough to affect a person's performance at work, school, and in everyday activities in a small percentage of women. When severe pelvic pain and bleeding suddenly occur or worsen during a cycle, this could be due to ectopic pregnancy and spontaneous abortion. This is checked by using a pregnancy test, ideally as soon as unusual pain begins, because ectopic pregnancies can be life‑threatening. The most common treatment for menstrual cramps are non-steroidal anti-inflammatory drugs (NSAIDs). NSAIDs can be used to reduce moderate to severe pain, and all appear similar. About 1 in 5 women do not respond to NSAIDs and require alternative therapy, such as simple analgesics or heat pads. Other medications for pain management include aspirin or paracetamol and combined oral contraceptives. Although combined oral contraceptives may be used, there is insufficient evidence for the efficacy of intrauterine progestogens. One review found tentative evidence that acupuncture may be useful, at least in the short term. Another review found insufficient evidence to determine an effect. Interactions with other conditions Known interactions between the menstrual cycle and certain health conditions include: Some women with neurological conditions experience increased activity of their conditions at about the same time during each menstrual cycle. For example, drops in estrogen levels may trigger migraines, especially when the woman who has migraines is also taking the birth control pill. Many women with epilepsy have more seizures in a pattern linked to the menstrual cycle; this is called "catamenial epilepsy". Different patterns seem to exist (such as seizures coinciding with the time of menstruation, or coinciding with the time of ovulation), and the frequency with which they occur has not been firmly established. Research indicates that women have a significantly higher likelihood of anterior cruciate ligament injuries in the pre-ovulatory stage, than post-ovulatory stage. Sexual activity Sexual feelings and behaviors change during the menstrual cycle. Before and during ovulation, high levels of estrogen and androgens result in women having a relatively increased interest in sexual activity, and relatively lower interest directly prior to and during menstruation. Unlike other mammals, women may show interest in sexual activity across all days of the menstrual cycle, regardless of fertility. There is no reliable scientific evidence that would advise against sexual intercourse during menstruation based on medical grounds. Fertility aspects Peak fertility (the time with the highest likelihood of pregnancy resulting from sexual intercourse) occurs during just a few days of the cycle: usually two days before and two days after the ovulation date. This corresponds to the second and the beginning of the third week in a 28-day cycle. This fertile window varies from woman to woman, just as the ovulation date often varies from cycle to cycle for the same woman. A variety of methods have been developed to help individual women estimate the relatively fertile and the relatively infertile days in the cycle; these systems are called fertility awareness. Menstrual disorders Infrequent or irregular ovulation is called oligoovulation. The absence of ovulation is called anovulation. Normal menstrual flow can occur without ovulation preceding it: an anovulatory cycle. In some cycles, follicular development may start but not be completed; nevertheless, estrogens will be formed and stimulate the uterine lining. Anovulatory flow resulting from a very thick endometrium caused by prolonged, continued high estrogen levels is called estrogen breakthrough bleeding. Anovulatory bleeding triggered by a sudden drop in estrogen levels is called withdrawal bleeding. Anovulatory cycles commonly occur before menopause (perimenopause) and in women with polycystic ovary syndrome. Very little flow (less than 10 ml) is called hypomenorrhea. Regular cycles with intervals of 21 days or fewer are polymenorrhea; frequent but irregular menstruation is known as metrorrhagia. Sudden heavy flows or amounts greater than 80 ml are termed menorrhagia. Heavy menstruation that occurs frequently and irregularly is menometrorrhagia. The term for cycles with intervals exceeding 35 days is oligomenorrhea. Amenorrhea refers to more than three to six months without menses (while not being pregnant) during a woman's reproductive years. The term for painful periods is dysmenorrhea. There is a wide spectrum of differences in how women experience menstruation. There are several ways that someone's menstrual cycle can differ from the norm: Extreme psychological stress can also result in periods stopping. More severe symptoms of anxiety or depression may be signs of premenstrual dysphoric disorder (PMDD) which is a depressive disorder. Dysfunctional uterine bleeding is a hormonally caused bleeding abnormality. Dysfunctional uterine bleeding typically occurs in premenopausal women who do not ovulate normally (i.e. are anovulatory). All these bleeding abnormalities need medical attention; they may indicate hormone imbalances, uterine fibroids, or other problems. As pregnant women may bleed, a pregnancy test forms part of the evaluation of abnormal bleeding. Women who had undergone female genital mutilation (particularly type III- infibulation) a practice common in parts of Africa, may experience menstrual problems, such as slow and painful menstruation, that is caused by the near-complete sealing off of the vagina. Dysmenorrhea Menstrual hygiene management Menstrual products (also called "feminine hygiene" products) are made to absorb or catch menstrual blood. A number of different products are available – some are disposable, some are reusable. Where women can afford it, items used to absorb or catch menses are usually commercially manufactured products. Menstruating women manage menstruation primarily by wearing menstrual products such as tampons, napkins or menstrual cups to catch the menstrual blood. The main disposable products (commercially manufactured) include: Sanitary napkins (also called sanitary towels or pads) – Rectangular pieces of material worn attached to the underwear to absorb menstrual flow, often with an adhesive backing to hold the pad in place. Disposable pads may contain wood pulp or gel products, usually with a plastic lining and bleached. Tampons – Disposable cylinders of treated rayon/cotton blends or all-cotton fleece, usually bleached, that are inserted into the vagina to absorb menstrual flow. The main reusable products include: Menstrual cups – A firm, flexible bell-shaped device worn inside the vagina to collect menstrual flow. Reusable cloth pads – Pads that are made of cotton (often organic), terrycloth, or flannel, and may be handsewn (from material or reused old clothes and towels) or storebought. Padded panties or period-proof underwear – Reusable cloth (usually cotton) underwear with extra absorbent layers sewn in to absorb flow. Due to poverty, some women cannot afford commercial feminine hygiene products. Instead, they use materials found in the environment or other improvised materials. "Period poverty" is a global issue affecting women and girls who do not have access to safe, hygienic sanitary products. In addition, solid waste disposal systems in developing countries are often lacking, which means women have no proper place to dispose used products, such as pads. Inappropriate disposal of used materials also creates pressures on sanitation systems as menstrual hygiene products can create blockages of toilets, pipes and sewers. In the UK research has shown that for women allotment growers, access to sanitation for menstrual hygiene management is limited. Menstrual suppression Due to hormonal contraception Menstruation can be delayed by the use of progesterone or progestins. For this purpose, oral administration of progesterone or progestin during cycle day 20 has been found to effectively delay menstruation for at least 20 days, with menstruation starting after 2–3 days have passed since discontinuing the regimen. Hormonal contraception affects the frequency, duration, severity, volume, and regularity of menstruation and menstrual symptoms. The most common form of hormonal contraception is the combined birth control pill, which contains both estrogen and progestogen. Although the primary function of the pill is to prevent pregnancy, it may be used to improve some menstrual symptoms and syndromes which affect menstruation, such as polycystic ovary syndrome (PCOS), endometriosis, adenomyosis, amenorrhea, menstrual cramps, menstrual migraines, menorrhagia (excessive menstrual bleeding), menstruation-related or fibroid-related anemia and dysmenorrhea (painful menstruation) by creating regularity in menstrual cycles and reducing overall menstrual flow. Using the combined birth control pill, it is also possible for a woman to delay or eliminate menstrual periods, a practice called menstrual suppression. Some women do this simply for convenience in the short-term, while others prefer to eliminate periods altogether when possible. This can be done either by skipping the placebo pills, or using an extended cycle combined oral contraceptive pill, which were first marketed in the U.S. in the early 2000s. This continuous administration of active pills without the placebo can lead to the achievement of amenorrhea in 80% of users within 1 year of use. Due to breastfeeding Breastfeeding causes negative feedback to occur on pulse secretion of gonadotropin-releasing hormone (GnRH) and luteinizing hormone (LH). Depending on the strength of the negative feedback, breastfeeding women may experience complete suppression of follicular development, follicular development but no ovulation, or normal menstrual cycles may resume. Suppression of ovulation is more likely when suckling occurs more frequently. The production of prolactin in response to suckling is important to maintaining lactational amenorrhea. On average, women who are fully breastfeeding whose infants suckle frequently experience a return of menstruation at fourteen and a half months postpartum. There is a wide range of response among individual breastfeeding women, however, with some experiencing return of menstruation at two months and others remaining amenorrheic for up to 42 months postpartum. Society and culture Etymology and terminology The word menstruation is etymologically related to moon. The terms menstruation and menses are derived from the Latin , which in turn relates to the ancient Greek and to the roots of the English words month and moon. Some organizations have begun to use the term "menstruator" instead of "menstruating women", a term that has been in use since at least 2010. Menstruator is used by activists and scholars in order to "express solidarity with women who do not menstruate, transgender men who do, and intersexual and genderqueer individuals". The term can be contentious between different schools of feminist thought; however, the majority of feminist scholars consider the term to correctly reflect the reality that people of different genders menstruate. The term "people who menstruate" is also used. Traditions, taboos and education Many religions have menstruation-related traditions, for example: Islam prohibits sexual contact with women during menstruation in the 2nd chapter of the Quran. Some scholars argue that menstruating women are in a state in which they are unable to maintain wudhu, and are therefore prohibited from touching the Arabic version of the Qur'an. In Judaism, a woman during menstruation is called Niddah and may be banned from certain actions. For example, the Jewish Torah prohibits sexual intercourse with a menstruating woman. In Hinduism, menstruating women are traditionally considered ritually impure and given rules to follow. In Zoroastrianism, if a woman’s menses did not stop after nine days, it was considered the work of the daēvas. Menstruation education is frequently taught in combination with sex education at school in Western countries, although girls may prefer their mothers to be the primary source of information about menstruation and puberty. Information about menstruation is often shared among friends and peers, which may promote a more positive outlook on puberty. The quality of menstrual education in a society determines the accuracy of people's understanding of the process. In many Western countries where menstruation is a taboo subject, girls tend to conceal the fact that they may be menstruating and struggle to ensure that they give no sign of menstruation. Effective educational programs are essential to providing children and adolescents with clear and accurate information about menstruation. Schools can be an appropriate place for menstrual education to take place. Programs led by peers or third-party agencies are another option. Low-income girls are less likely to receive proper sex education on puberty, leading to a decreased understanding of why menstruation occurs and the associated physiological changes that take place. This has been shown to cause the development of a negative attitude towards menstruation. Seclusion during menstruation In some cultures, women were isolated during menstruation due to menstrual taboos. This is because they are seen as unclean, dangerous, or bringing bad luck to those who encounter them. These practices are common in parts of South Asia including India. A 1983 report found women refraining from household chore during this period in India. Chhaupadi is a social practice that occurs in the western part of Nepal for Hindu women, which prohibits a woman from participating in everyday activities during menstruation. Women are considered impure during this time and are kept out of the house and have to live in a shed. Although chhaupadi was outlawed by the Supreme Court of Nepal in 2005, the tradition is slow to change. Women and girls in cultures which practice such seclusion are often confined to menstruation huts, which are places of isolation used by cultures with strong menstrual taboos. The practice has recently come under fire due to related fatalities. Nepal criminalized the practice in 2017 after deaths were reported after the elongated isolation periods, but "the practice of isolating menstruating women and girls continues." Not all cultures villainize menstruation, the Beng people of West Africa consider menstrual blood as sacred and recognize its significance in reproduction. Beliefs around synchrony Effects of the moon Even though the average length of the human menstrual cycle is similar to that of the lunar cycle, in modern humans there is no relation between the two. The relationship is believed to be a coincidence. Light exposure does not appear to affect the menstrual cycle in humans. A meta-analysis of studies from 1996 showed no correlation between the human menstrual cycle and the lunar cycle, nor did data analyzed by period-tracking app Clue, submitted by 1.5million women, of 7.5million menstrual cycles; however, the lunar cycle and the average menstrual cycle were found to be basically equal in length. Cohabitation Beginning in 1971, some research suggested that menstrual cycles of cohabiting women became synchronized (menstrual synchrony). Subsequent research has called this hypothesis into question. A 2013 review concluded that menstrual synchrony likely does not exist. Work Some countries, mainly in Asia, have menstrual leave to provide women with either paid or unpaid leave of absence from their employment while they are menstruating. Countries with policies include Japan, Taiwan, Indonesia, and South Korea. The practice is controversial due to concerns that it bolsters the perception of women as weak, inefficient workers, as well as concerns that it is unfair to men, and that it furthers gender stereotypes and the medicalization of menstruation. Other mammals Most female mammals have an estrous cycle, but not all have a menstrual cycle that results in menstruation. Menstruation in mammals occurs in some close evolutionary relatives such as chimpanzees.	Biology and health sciences	Human reproduction	Biology
38238	https://en.wikipedia.org/wiki/Hepatitis	Hepatitis	Hepatitis is inflammation of the liver tissue. Some people or animals with hepatitis have no symptoms, whereas others develop yellow discoloration of the skin and whites of the eyes (jaundice), poor appetite, vomiting, tiredness, abdominal pain, and diarrhea. Hepatitis is acute if it resolves within six months, and chronic if it lasts longer than six months. Acute hepatitis can resolve on its own, progress to chronic hepatitis, or (rarely) result in acute liver failure. Chronic hepatitis may progress to scarring of the liver (cirrhosis), liver failure, and liver cancer. Hepatitis is most commonly caused by the virus hepatovirus A, B, C, D, and E. Other viruses can also cause liver inflammation, including cytomegalovirus, Epstein–Barr virus, and yellow fever virus. Other common causes of hepatitis include heavy alcohol use, certain medications, toxins, other infections, autoimmune diseases, and non-alcoholic steatohepatitis (NASH). Hepatitis A and E are mainly spread by contaminated food and water. Hepatitis B is mainly sexually transmitted, but may also be passed from mother to baby during pregnancy or childbirth and spread through infected blood. Hepatitis C is commonly spread through infected blood such as may occur during needle sharing by intravenous drug users. Hepatitis D can only infect people already infected with hepatitis B. Hepatitis A, B, and D are preventable with immunization. Medications may be used to treat chronic viral hepatitis. Antiviral medications are recommended in all with chronic hepatitis C, except those with conditions that limit their life expectancy. There is no specific treatment for NASH; physical activity, a healthy diet, and weight loss are recommended. Autoimmune hepatitis may be treated with medications to suppress the immune system. A liver transplant may be an option in both acute and chronic liver failure. Worldwide in 2015, hepatitis A occurred in about 114 million people, chronic hepatitis B affected about 343 million people and chronic hepatitis C about 142 million people. In the United States, NASH affects about 11 million people and alcoholic hepatitis affects about 5 million people. Hepatitis results in more than a million deaths a year, most of which occur indirectly from liver scarring or liver cancer. In the United States, hepatitis A is estimated to occur in about 2,500 people a year and results in about 75 deaths. The word is derived from the Greek hêpar (), meaning "liver", and -itis (), meaning "inflammation". Signs and symptoms Hepatitis has a broad spectrum of presentations that range from a complete lack of symptoms to severe liver failure. The acute form of hepatitis, generally caused by viral infection, is characterized by constitutional symptoms that are typically self-limiting. Chronic hepatitis presents similarly, but can manifest signs and symptoms specific to liver dysfunction with long-standing inflammation and damage to the organ. Acute hepatitis Acute viral hepatitis follows three distinct phases: The initial prodromal phase (preceding symptoms) involves non-specific and flu-like symptoms common to many acute viral infections. These include fatigue, nausea, vomiting, poor appetite, joint pain, and headaches. Fever, when present, is most common in cases of hepatitis A and E. Late in this phase, people can experience liver-specific symptoms, including choluria (dark urine) and clay-colored stools. Yellowing of the skin and whites of the eyes follow the prodrome after about 1–2 weeks and can last for up to 4 weeks. The non-specific symptoms seen in the prodromal typically resolve by this time, but people will develop an enlarged liver and right upper abdominal pain or discomfort. 10–20% of people will also experience an enlarged spleen, while some people will also experience a mild unintentional weight loss. The recovery phase is characterized by resolution of the clinical symptoms of hepatitis with persistent elevations in liver lab values and potentially a persistently enlarged liver. All cases of hepatitis A and E are expected to fully resolve after 1–2 months. Most hepatitis B cases are also self-limiting and will resolve in 3–4 months. Few cases of hepatitis C will resolve completely. Both drug-induced hepatitis and autoimmune hepatitis can present very similarly to acute viral hepatitis, with slight variations in symptoms depending on the cause. Cases of drug-induced hepatitis can manifest with systemic signs of an allergic reaction including rash, fever, serositis (inflammation of membranes lining certain organs), elevated eosinophils (a type of white blood cell), and suppression of bone marrow activity. Fulminant hepatitis Fulminant hepatitis, or massive hepatic cell death, is a rare and life-threatening complication of acute hepatitis that can occur in cases of hepatitis B, D, and E, in addition to drug-induced and autoimmune hepatitis. The complication more frequently occurs in instances of hepatitis B and D co-infection at a rate of 2–20% and in pregnant women with hepatitis E at rate of 15–20% of cases. In addition to the signs of acute hepatitis, people can also demonstrate signs of coagulopathy (abnormal coagulation studies with easy bruising and bleeding) and encephalopathy (confusion, disorientation, and sleepiness). Mortality due to fulminant hepatitis is typically the result of various complications including cerebral edema, gastrointestinal bleeding, sepsis, respiratory failure, or kidney failure. Chronic hepatitis Acute cases of hepatitis are seen to be resolved well within a six-month period. When hepatitis is continued for more than six months it is termed chronic hepatitis. Chronic hepatitis is often asymptomatic early in its course and is detected only by liver laboratory studies for screening purposes or to evaluate non-specific symptoms. As the inflammation progresses, patients can develop constitutional symptoms similar to acute hepatitis, including fatigue, nausea, vomiting, poor appetite, and joint pain. Jaundice can occur as well, but much later in the disease process and is typically a sign of advanced disease. Chronic hepatitis interferes with hormonal functions of the liver which can result in acne, hirsutism (abnormal hair growth), and amenorrhea (lack of menstrual period) in women. Extensive damage and scarring of the liver over time defines cirrhosis, a condition in which the liver's ability to function is permanently impeded. This results in jaundice, weight loss, coagulopathy, ascites (abdominal fluid collection), and peripheral edema (leg swelling). Cirrhosis can lead to other life-threatening complications such as hepatic encephalopathy, esophageal varices, hepatorenal syndrome, and liver cancer. Causes Causes of hepatitis can be divided into the following major categories: infectious, metabolic, ischemic, autoimmune, genetic, and other. Infectious agents include viruses, bacteria, and parasites. Metabolic causes include prescription medications, toxins (most notably alcohol), and non-alcoholic fatty liver disease. Autoimmune and genetic causes of hepatitis involve genetic predispositions and tend to affect characteristic populations. Infectious Viral hepatitis Viral hepatitis is the most common type of hepatitis worldwide, especially in Asia and Africa. Viral hepatitis is caused by five different viruses (hepatitis A, B, C, D, and E). Hepatitis A and hepatitis E behave similarly: they are both transmitted by the fecal–oral route, are more common in developing countries, and are self-limiting illnesses that do not lead to chronic hepatitis. Hepatitis B, hepatitis C, and hepatitis D are transmitted when blood or mucous membranes are exposed to infected blood and body fluids, such as semen and vaginal secretions. Viral particles have also been found in saliva and breastmilk. Kissing, sharing utensils, and breastfeeding do not lead to transmission unless these fluids are introduced into open sores or cuts. Many families who do not have safe drinking water or live in unhygienic homes have contracted hepatitis because saliva and blood droplets are often carried through the water and blood-borne illnesses spread quickly in unsanitary settings. Hepatitis B and C can present either acutely or chronically. Hepatitis D is a defective virus that requires hepatitis B to replicate and is only found with hepatitis B co-infection. In adults, hepatitis B infection is most commonly self-limiting, with less than 5% progressing to chronic state, and 20 to 30% of those chronically infected developing cirrhosis or liver cancer. Infection in infants and children frequently leads to chronic infection. Unlike hepatitis B, most cases of hepatitis C lead to chronic infection. Hepatitis C is the second most common cause of cirrhosis in the US (second to alcoholic hepatitis). In the 1970s and 1980s, blood transfusions were a major factor in spreading hepatitis C virus. Since widespread screening of blood products for hepatitis C began in 1992, the risk of acquiring hepatitis C from a blood transfusion has decreased from approximately 10% in the 1970s to 1 in 2 million currently. Parasitic hepatitis Parasites can also infect the liver and activate the immune response, resulting in symptoms of acute hepatitis with increased serum IgE (though chronic hepatitis is possible with chronic infections). Of the protozoans, Trypanosoma cruzi, Leishmania species, and the malaria-causing Plasmodium species all can cause liver inflammation. Another protozoan, Entamoeba histolytica, causes hepatitis with distinct liver abscesses. Of the worms, the cestode Echinococcus granulosus, also known as the dog tapeworm, infects the liver and forms characteristic hepatic hydatid cysts. The liver flukes Fasciola hepatica and Clonorchis sinensis live in the bile ducts and cause progressive hepatitis and liver fibrosis. Bacterial hepatitis Bacterial infection of the liver commonly results in pyogenic liver abscesses, acute hepatitis, or granulomatous (or chronic) liver disease. Pyogenic abscesses commonly involve enteric bacteria such as Escherichia coli and Klebsiella pneumoniae and are composed of multiple bacteria up to 50% of the time. Acute hepatitis is caused by Neisseria meningitidis, Neisseria gonorrhoeae, Bartonella henselae, Borrelia burgdorferi, salmonella species, brucella species and campylobacter species. Chronic or granulomatous hepatitis is seen with infection from mycobacteria species, Tropheryma whipplei, Treponema pallidum, Coxiella burnetii, and rickettsia species. Metabolic Alcoholic hepatitis Excessive alcohol consumption is a significant cause of hepatitis and is the most common cause of cirrhosis in the U.S. Alcoholic hepatitis is within the spectrum of alcoholic liver disease. This ranges in order of severity and reversibility from alcoholic steatosis (least severe, most reversible), alcoholic hepatitis, cirrhosis, and liver cancer (most severe, least reversible). Hepatitis usually develops over years-long exposure to alcohol, occurring in 10 to 20% of alcoholics. The most important risk factors for the development of alcoholic hepatitis are quantity and duration of alcohol intake. Long-term alcohol intake in excess of 80 grams of alcohol a day in men and 40 grams a day in women is associated with development of alcoholic hepatitis (1 beer or 4 ounces of wine is equivalent to 12g of alcohol). Alcoholic hepatitis can vary from asymptomatic hepatomegaly (enlarged liver) to symptoms of acute or chronic hepatitis to liver failure. Toxic and drug-induced hepatitis Many chemical agents, including medications, industrial toxins, and herbal and dietary supplements, can cause hepatitis. The spectrum of drug-induced liver injury varies from acute hepatitis to chronic hepatitis to acute liver failure. Toxins and medications can cause liver injury through a variety of mechanisms, including direct cell damage, disruption of cell metabolism, and causing structural changes. Some drugs such as paracetamol exhibit predictable dose-dependent liver damage while others such as isoniazid cause idiosyncratic and unpredictable reactions that vary by person. There are wide variations in the mechanisms of liver injury and latency period from exposure to development of clinical illness. Many types of drugs can cause liver injury, including the analgesic paracetamol; antibiotics such as isoniazid, nitrofurantoin, amoxicillin-clavulanate, erythromycin, and trimethoprim-sulfamethoxazole; anticonvulsants such as valproate and phenytoin; cholesterol-lowering statins; steroids such as oral contraceptives and anabolic steroids; and highly active anti-retroviral therapy used in the treatment of HIV/AIDS. Of these, amoxicillin-clavulanate is the most common cause of drug-induced liver injury, and paracetamol toxicity the most common cause of acute liver failure in the United States and Europe. Herbal remedies and dietary supplements are another important cause of hepatitis; these are the most common causes of drug-induced hepatitis in Korea. The United States–based Drug Induced Liver Injury Network linked more than 16% of cases of hepatotoxicity to herbal and dietary supplements. In the United States, herbal and dietary supplements – unlike pharmaceutical drugs – are unregulated by the Food and Drug Administration. The National Institutes of Health maintains the LiverTox database for consumers to track all known prescription and non-prescription compounds associated with liver injury. Exposure to other hepatotoxins can occur accidentally or intentionally through ingestion, inhalation, and skin absorption. The industrial toxin carbon tetrachloride and the wild mushroom Amanita phalloides are other known hepatotoxins. Non-alcoholic fatty liver disease Non-alcoholic hepatitis is within the spectrum of non-alcoholic liver disease (NALD), which ranges in severity and reversibility from non-alcoholic fatty liver disease (NAFLD) to non-alcoholic steatohepatitis (NASH) to cirrhosis to liver cancer, similar to the spectrum of alcoholic liver disease. Non-alcoholic liver disease occurs in people with little or no history of alcohol use, and is instead strongly associated with metabolic syndrome, obesity, insulin resistance and diabetes, and hypertriglyceridemia. Over time, non-alcoholic fatty liver disease can progress to non-alcoholic steatohepatitis, which additionally involves liver cell death, liver inflammation and possible fibrosis. Factors accelerating progression from NAFLD to NASH are obesity, older age, non-African American ethnicity, female gender, diabetes mellitus, hypertension, higher ALT or AST level, higher AST/ALT ratio, low platelet count, and an ultrasound steatosis score. In the early stages (as with NAFLD and early NASH), most patients are asymptomatic or have mild right upper quadrant pain, and diagnosis is suspected on the basis of abnormal liver function tests. As the disease progresses, symptoms typical of chronic hepatitis may develop. While imaging can show fatty liver, only liver biopsy can demonstrate inflammation and fibrosis characteristic of NASH. 9 to 25% of patients with NASH develop cirrhosis. NASH is recognized as the third most common cause of liver disease in the United States. Autoimmune Autoimmune hepatitis is a chronic disease caused by an abnormal immune response against liver cells. The disease is thought to have a genetic predisposition as it is associated with certain human leukocyte antigens involved in the immune response. As in other autoimmune diseases, circulating auto-antibodies may be present and are helpful in diagnosis. Auto-antibodies found in patients with autoimmune hepatitis include the sensitive but less specific anti-nuclear antibody (ANA), smooth muscle antibody (SMA), and atypical perinuclear antineutrophil cytoplasmic antibody (p-ANCA). Other autoantibodies that are less common but more specific to autoimmune hepatitis are the antibodies against liver kidney microsome 1 (LKM1) and soluble liver antigen (SLA). Autoimmune hepatitis can also be triggered by drugs (such as nitrofurantoin, hydralazine, and methyldopa), after liver transplant, or by viruses (such as hepatitis A, Epstein-Barr virus, or measles). Autoimmune hepatitis can present anywhere within the spectrum from asymptomatic to acute or chronic hepatitis to fulminant liver failure. Patients are asymptomatic 25–34% of the time, and the diagnosis is suspected on the basis of abnormal liver function tests. Some studies show between 25% and 75% of cases present with signs and symptoms of acute hepatitis. As with other autoimmune diseases, autoimmune hepatitis usually affects young females (though it can affect patients of either sex of any age), and patients can exhibit classic signs and symptoms of autoimmunity such as fatigue, anemia, anorexia, amenorrhea, acne, arthritis, pleurisy, thyroiditis, ulcerative colitis, nephritis, and maculopapular rash. Autoimmune hepatitis increases the risk for cirrhosis, and the risk for liver cancer is increased by about 1% for each year of the disease. Many people with autoimmune hepatitis have other autoimmune diseases. Autoimmune hepatitis is distinct from the other autoimmune diseases of the liver, primary biliary cirrhosis and primary sclerosing cholangitis, both of which can also lead to scarring, fibrosis, and cirrhosis of the liver. Genetic Genetic causes of hepatitis include alpha-1-antitrypsin deficiency, hemochromatosis, and Wilson's disease. In alpha-1-antitrypsin deficiency, a co-dominant mutation in the gene for alpha-1-antitrypsin results in the abnormal accumulation of the mutant AAT protein within liver cells, leading to liver disease. Hemochromatosis and Wilson's disease are both autosomal recessive diseases involving abnormal storage of minerals. In hemochromatosis, excess amounts of iron accumulate in multiple body sites, including the liver, which can lead to cirrhosis. In Wilson's disease, excess amounts of copper accumulate in the liver and brain, causing cirrhosis and dementia. When the liver is involved, alpha-1-antitrypsin deficiency and Wilson's disease tend to present as hepatitis in the neonatal period or in childhood. Hemochromatosis typically presents in adulthood, with the onset of clinical disease usually after age 50. Ischemic hepatitis Ischemic hepatitis (also known as shock liver) results from reduced blood flow to the liver as in shock, heart failure, or vascular insufficiency. The condition is most often associated with heart failure but can also be caused by shock or sepsis. Blood testing of a person with ischemic hepatitis will show very high levels of transaminase enzymes (AST and ALT). The condition usually resolves if the underlying cause is treated successfully. Ischemic hepatitis rarely causes permanent liver damage. Other Hepatitis can also occur in neonates and is attributable to a variety of causes, some of which are not typically seen in adults. Congenital or perinatal infection with the hepatitis viruses, toxoplasma, rubella, cytomegalovirus, and syphilis can cause neonatal hepatitis. Structural abnormalities such as biliary atresia and choledochal cysts can lead to cholestatic liver injury leading to neonatal hepatitis. Metabolic diseases such as glycogen storage disorders and lysosomal storage disorders are also implicated. Neonatal hepatitis can be idiopathic, and in such cases, biopsy often shows large multinucleated cells in the liver tissue. This disease is termed giant cell hepatitis and may be associated with viral infection, autoimmune disorders, and drug toxicity. Mechanism The specific mechanism varies and depends on the underlying cause of the hepatitis. Generally, there is an initial insult that causes liver injury and activation of an inflammatory response, which can become chronic, leading to progressive fibrosis and cirrhosis. Viral hepatitis The pathway by which hepatic viruses cause viral hepatitis is best understood in the case of hepatitis B and C. The viruses do not directly activate apoptosis (cell death). Rather, infection of liver cells activates the innate and adaptive arms of the immune system leading to an inflammatory response which causes cellular damage and death, including viral-induced apoptosis via the induction of the death receptor-mediated signaling pathway. Depending on the strength of the immune response, the types of immune cells involved and the ability of the virus to evade the body's defense, infection can either lead to clearance (acute disease) or persistence (chronic disease) of the virus. The chronic presence of the virus within liver cells results in multiple waves of inflammation, injury and wound healing that over time lead to scarring or fibrosis and culminate in hepatocellular carcinoma. People with impaired immune response are at greater risk of developing chronic infection. Natural killer cells are the primary drivers of the initial innate response and create a cytokine environment that results in the recruitment of CD4 T-helper and CD8 cytotoxic T-cells. Type I interferons are the cytokines that drive the antiviral response. In chronic Hepatitis B and C, natural killer cell function is impaired. Steatohepatitis Steatohepatitis is seen in both alcoholic and non-alcoholic liver disease and is the culmination of a cascade of events that began with injury. In the case of non-alcoholic steatohepatitis, this cascade is initiated by changes in metabolism associated with obesity, insulin resistance, and lipid dysregulation. In alcoholic hepatitis, chronic excess alcohol use is the culprit. Though the inciting event may differ, the progression of events is similar and begins with accumulation of free fatty acids (FFA) and their breakdown products in the liver cells in a process called steatosis. This initially reversible process overwhelms the hepatocyte's ability to maintain lipid homeostasis leading to a toxic effect as fat molecules accumulate and are broken down in the setting of an oxidative stress response. Over time, this abnormal lipid deposition triggers the immune system via toll-like receptor 4 (TLR4) resulting in the production of inflammatory cytokines such as TNF that cause liver cell injury and death. These events mark the transition to steatohepatitis and in the setting of chronic injury, fibrosis eventually develops setting up events that lead to cirrhosis and hepatocellular carcinoma. Microscopically, changes that can be seen include steatosis with large and swollen hepatocytes (ballooning), evidence of cellular injury and cell death (apoptosis, necrosis), evidence of inflammation in particular in zone 3 of the liver, variable degrees of fibrosis and Mallory bodies. Diagnosis Diagnosis of hepatitis is made on the basis of some or all of the following: a person's signs and symptoms, medical history including sexual and substance use history, blood tests, imaging, and liver biopsy. In general, for viral hepatitis and other acute causes of hepatitis, the person's blood tests and clinical picture are sufficient for diagnosis. For other causes of hepatitis, especially chronic causes, blood tests may not be useful. In this case, liver biopsy is the gold standard for establishing the diagnosis: histopathologic analysis is able to reveal the precise extent and pattern of inflammation and fibrosis. Biopsy is typically not the initial diagnostic test because it is invasive and is associated with a small but significant risk of bleeding that is increased in people with liver injury and cirrhosis. Blood testing includes liver enzymes, serology (i.e. for autoantibodies), nucleic acid testing (i.e. for hepatitis virus DNA/RNA), blood chemistry, and complete blood count. Characteristic patterns of liver enzyme abnormalities can point to certain causes or stages of hepatitis. Generally, AST and ALT are elevated in most cases of hepatitis regardless of whether the person shows any symptoms. The degree of elevation (i.e. levels in the hundreds vs. in the thousands), the predominance for AST vs. ALT elevation, and the ratio between AST and ALT are informative of the diagnosis. Ultrasound, CT, and MRI can all identify steatosis (fatty changes) of the liver tissue and nodularity of the liver surface suggestive of cirrhosis. CT and especially MRI are able to provide a higher level of detail, allowing visualization and characterize such structures as vessels and tumors within the liver. Unlike steatosis and cirrhosis, no imaging test is able to detect liver inflammation (i.e. hepatitis) or fibrosis. Liver biopsy is the only definitive diagnostic test that is able to assess inflammation and fibrosis of the liver. Viral hepatitis Viral hepatitis is primarily diagnosed through blood tests for levels of viral antigens (such as the hepatitis B surface or core antigen), anti-viral antibodies (such as the anti-hepatitis B surface antibody or anti-hepatitis A antibody), or viral DNA/RNA. In early infection (i.e. within 1 week), IgM antibodies are found in the blood. In late infection and after recovery, IgG antibodies are present and remain in the body for up to years. Therefore, when a patient is positive for IgG antibody but negative for IgM antibody, he is considered immune from the virus via either prior infection and recovery or prior vaccination. In the case of hepatitis B, blood tests exist for multiple virus antigens (which are different components of the virion particle) and antibodies. The combination of antigen and antibody positivity can provide information about the stage of infection (acute or chronic), the degree of viral replication, and the infectivity of the virus. Alcoholic versus non-alcoholic The most apparent distinguishing factor between alcoholic steatohepatitis (ASH) and nonalcoholic steatohepatitis (NASH) is a history of excessive alcohol use. Thus, in patients who have no or negligible alcohol use, the diagnosis is unlikely to be alcoholic hepatitis. In those who drink alcohol, the diagnosis may just as likely be alcoholic or nonalcoholic hepatitis especially if there is concurrent obesity, diabetes, and metabolic syndrome. In this case, alcoholic and nonalcoholic hepatitis can be distinguished by the pattern of liver enzyme abnormalities; specifically, in alcoholic steatohepatitis AST>ALT with ratio of AST:ALT>2:1 while in nonalcoholic steatohepatitis ALT>AST with ratio of ALT:AST>1.5:1. Liver biopsies show identical findings in patients with ASH and NASH, specifically, the presence of polymorphonuclear infiltration, hepatocyte necrosis and apoptosis in the form of ballooning degeneration, Mallory bodies, and fibrosis around veins and sinuses. Virus screening The purpose of screening for viral hepatitis is to identify people infected with the disease as early as possible, even before symptoms and transaminase elevations may be present. This allows for early treatment, which can both prevent disease progression and decrease the likelihood of transmission to others. Hepatitis A Hepatitis A causes an acute illness that does not progress to chronic liver disease. Therefore, the role of screening is to assess immune status in people who are at high risk of contracting the virus, as well as in people with known liver disease for whom hepatitis A infection could lead to liver failure. People in these groups who are not already immune can receive the hepatitis A vaccine. Those at high risk and in need of screening include: People in close contact (either living with or having sexual contact) with someone who has hepatitis A People traveling to an area with endemic hepatitis A People who do not have access to clean water People who use illicit drugs People with liver disease People with poor sanitary habits such as not washing hands after using the restroom or changing diapers The presence of anti-hepatitis A IgG in the blood indicates past infection with the virus or prior vaccination. Hepatitis B The CDC, WHO, USPSTF, and ACOG recommend routine hepatitis B screening for certain high-risk populations. Specifically, these populations include people who are: Beginning immunosuppressive or cytotoxic therapy Blood, organ, or tissue donors Born in countries where the prevalence of hepatitis B is high (defined as ≥2% of the population), whether or not they have been vaccinated Born in the United States whose parents are from countries where the prevalence of hepatitis B is very high (defined as ≥8% of the population), and who were not vaccinated Found to have elevated liver enzymes without a known cause HIV positive In close contact with (i.e. live or have sex with) people known to have hepatitis B Incarcerated Intravenous drug users Men who have sex with men On hemodialysis Pregnant Screening consists of a blood test that detects hepatitis B surface antigen (HBsAg). If HBsAg is present, a second test – usually done on the same blood sample – that detects the antibody for the hepatitis B core antigen (anti-HBcAg) can differentiate between acute and chronic infection. People who are high-risk whose blood tests negative for HBsAg can receive the hepatitis B vaccine to prevent future infection. Hepatitis C The CDC, WHO, USPSTF, AASLD, and ACOG recommend screening people at high risk for hepatitis C infection. These populations include people who are: Adults in the United States born between 1945 and 1965 Blood or organ donors. Born to HCV-positive mothers HIV-positive Incarcerated, or who have been in the past Intranasal illicit drug users Intravenous drug users (past or current) Men who have sex with men On long-term hemodialysis, or who have been in the past Pregnant, and engaging in high-risk behaviors Recipients of blood products or organs prior to 1992 in the United States Recipients of tattoos in an "unregulated setting" Sex workers Workers in a healthcare setting who have had a needlestick injury For people in the groups above whose exposure is ongoing, screening should be periodic, though there is no set optimal screening interval. The AASLD recommends screening men who have sex with men who are HIV-positive annually. People born in the US between 1945 and 1965 should be screened once (unless they have other exposure risks). Screening consists of a blood test that detects anti-hepatitis C virus antibody. If anti-hepatitis C virus antibody is present, a confirmatory test to detect HCV RNA indicates chronic disease. Hepatitis D The CDC, WHO, USPSTF, AASLD, and ACOG recommend screening people at high risk for hepatitis D infection. These populations include people who are: Blood or organ donors. Incarcerated, or who have been in the past Intranasal illicit drug users Intravenous drug users (past or current) Sex workers Workers in a healthcare setting who have had a needlestick injury Hepatitis D is extremely rare. Symptoms include chronic diarrhea, anal and intestinal blisters, purple urine, and burnt popcorn scented breath. Screening consists of a blood test that detects the anti-hepitits D virus antibbody. If anti-hepitits D virus antibody is present, a confirmatory test to detect HDV RNA DNA indicates chronic disease. Prevention Vaccines Hepatitis A The CDC recommends the hepatitis A vaccine for all children beginning at age one, as well as for those who have not been previously immunized and are at high risk for contracting the disease. For children 12 months of age or older, the vaccination is given as a shot into the muscle in two doses 6–18 months apart and should be started before the age 24 months. The dosing is slightly different for adults depending on the type of the vaccine. If the vaccine is for hepatitis A only, two doses are given 6–18 months apart depending on the manufacturer. If the vaccine is combined hepatitis A and hepatitis B, up to 4 doses may be required. Hepatitis B The CDC recommends the routine vaccination of all children under the age of 19 with the hepatitis B vaccine. They also recommend it for those who desire it or are at high risk. Routine vaccination for hepatitis B starts with the first dose administered as a shot into the muscle before the newborn is discharged from the hospital. An additional two doses should be administered before the child is 18 months. For babies born to a mother with hepatitis B surface antigen positivity, the first dose is unique – in addition to the vaccine, the hepatitis immune globulin should also be administered, both within 12 hours of birth. These newborns should also be regularly tested for infection for at least the first year of life. There is also a combination formulation that includes both hepatitis A and B vaccines. Other There are currently no vaccines available in the United States for hepatitis C or E. In 2015, a group in China published an article regarding the development of a vaccine for hepatitis E. As of March 2016, the United States government was in the process of recruiting participants for the phase IV trial of the hepatitis E vaccine. Behavioral changes Hepatitis A Because hepatitis A is transmitted primarily through the oral-fecal route, the mainstay of prevention aside from vaccination is good hygiene, access to clean water and proper handling of sewage. Hepatitis B and C As hepatitis B and C are transmitted through blood and multiple bodily fluids, prevention is aimed at screening blood prior to transfusion, abstaining from the use of injection drugs, safe needle and sharps practices in healthcare settings, and safe sex practices. Hepatitis D The hepatitis D virus requires that a person first be infected with hepatitis B virus, so prevention efforts should focus on limiting the spread of hepatitis B. In people who have chronic hepatitis B infection and are at risk for superinfection with the hepatitis D virus, the preventive strategies are the same as for hepatitis B. Hepatitis E Hepatitis E is spread primarily through the oral-fecal route but may also be spread by blood and from mother to fetus. The mainstay of hepatitis E prevention is similar to that for hepatitis A (namely, good hygiene and clean water practices). Alcoholic and metabolic hepatitis As excessive alcohol consumption can lead to hepatitis and cirrhosis, the following are maximal recommendations for alcohol consumption: Men – ≤ 4 drinks on any given day and ≤ 14 drinks per week Women – ≤ 3 drinks on any given day and ≤ 7 drinks per week To prevent MAFLD it is recommended to maintain a normal weight, eat a healthy diet, avoid added sugar, and exercise regularly. Successes Hepatitis A In the United States, universal immunization has led to a two-thirds decrease in hospital admissions and medical expenses due to hepatitis A. Hepatitis B In the United States new cases of hepatitis B decreased 75% from 1990 to 2004. The group that saw the greatest decrease was children and adolescents, likely reflecting the implementation of the 1999 guidelines. Hepatitis C Hepatitis C infections each year had been declining since the 1980s, but began to increase again in 2006. The data are unclear as to whether the decline can be attributed to needle exchange programmes. Alcoholic hepatitis Because people with alcoholic hepatitis may have no symptoms, it can be difficult to diagnose and the number of people with the disease is probably higher than many estimates. Programs such as Alcoholics Anonymous have been successful in decreasing death due to cirrhosis, but it is difficult to evaluate their success in decreasing the incidence of alcoholic hepatitis. Treatment The treatment of hepatitis varies according to the type, whether it is acute or chronic, and the severity of the disease. Activity: Many people with hepatitis prefer bed rest, though it is not necessary to avoid all physical activity while recovering. Diet: A high-calorie diet is recommended. Many people develop nausea and cannot tolerate food later in the day, so the bulk of intake may be concentrated in the earlier part of the day. In the acute phase of the disease, intravenous feeding may be needed if patients cannot tolerate food and have poor oral intake subsequent to nausea and vomiting. Drugs: People with hepatitis should avoid taking drugs metabolized by the liver. Glucocorticoids are not recommended as a treatment option for acute viral hepatitis and may even cause harm, such as development of chronic hepatitis. Precautions: Universal precautions should be observed. Isolation is usually not needed, except in cases of hepatitis A and E who have fecal incontinence, and in cases of hepatitis B and C who have uncontrolled bleeding. Hepatitis A Hepatitis A usually does not progress to a chronic state, and rarely requires hospitalization. Treatment is supportive and includes such measures as providing intravenous (IV) hydration and maintaining adequate nutrition. Rarely, people with the hepatitis A virus can rapidly develop liver failure, termed fulminant hepatic failure, especially the elderly and those who had a pre-existing liver disease, especially hepatitis C. Mortality risk factors include greater age and chronic hepatitis C. In these cases, more aggressive supportive therapy and liver transplant may be necessary. Hepatitis B Acute In healthy patients, 95–99% recover with no long-lasting effects, and antiviral treatment is not warranted. Age and comorbid conditions can result in a more prolonged and severe illness. Certain patients warrant hospitalization, especially those who present with clinical signs of ascites, peripheral edema, and hepatic encephalopathy, and laboratory signs of hypoglycemia, prolonged prothrombin time, low serum albumin, and very high serum bilirubin. In these rare, more severe acute cases, patients have been successfully treated with antiviral therapy similar to that used in cases of chronic hepatitis B, with nucleoside analogues such as entecavir or tenofovir. As there is a dearth of clinical trial data and the drugs used to treat are prone to developing resistance, experts recommend reserving treatment for severe acute cases, not mild to moderate. Chronic Chronic hepatitis B management aims to control viral replication, which is correlated with progression of disease. Seven drugs are approved in the United States: Adefovir dipivoxil, a nucleotide analogue, has been used to supplement lamivudine in patients who develop resistance, but is no longer recommended as first-line therapy. Entecavir is safe, well tolerated, less prone to developing resistance, and the most potent of the existing hepatitis B antivirals; it is thus a first-line treatment choice. It is not recommended for lamivudine-resistant patients or as monotherapy in patients who are HIV positive. Injectable interferon alpha was the first therapy approved for chronic hepatitis B. It has several side effects, most of which are reversible with removal of therapy, but it has been supplanted by newer treatments for this indication. These include long-acting interferon bound to polyethylene glycol (pegylated interferon) and the oral nucleoside analogues. Lamivudine was the first approved oral nucleoside analogue. While effective and potent, lamivudine has been replaced by newer, more potent treatments in the Western world and is no longer recommended as first-line treatment. It is still used in areas where newer agents either have not been approved or are too costly. Generally, the course of treatment is a minimum of one year with a minimum of six additional months of "consolidation therapy." Based on viral response, longer therapy may be required, and certain patients require indefinite long-term therapy. Due to a less robust response in Asian patients, consolidation therapy is recommended to be extended to at least a year. All patients should be monitored for viral reactivation, which if identified, requires restarting treatment. Lamivudine is generally safe and well tolerated. Many patients develop resistance, which is correlated with longer treatment duration. If this occurs, an additional antiviral is added. Lamivudine as a single treatment is contraindicated in patients coinfected with HIV, as resistance develops rapidly, but it can be used as part of a multidrug regimen. Pegylated interferon (PEG IFN) is dosed just once a week as a subcutaneous injection and is both more convenient and effective than standard interferon. Although it does not develop resistance as do many of the oral antivirals, it is poorly tolerated and requires close monitoring. PEG IFN is estimated to cost about $18,000 per year in the United States, compared to $2,500–8,700 for the oral medications. Its treatment duration is 48 weeks, unlike oral antivirals which require indefinite treatment for most patients (minimum one year). PEG IFN is not effective in patients with high levels of viral activity and cannot be used in immunosuppressed patients or those with cirrhosis. Telbivudine is effective but not recommended as first-line treatment; as compared to entecavir, it is both less potent and more resistance prone. Tenofovir is a nucleotide analogue and an antiretroviral drug that is also used to treat HIV infection. It is preferred to adefovir both in lamivudine-resistant patients and as initial treatment since it is both more potent and less likely to develop resistance. First-line treatments currently used include PEG IFN, entecavir, and tenofovir, subject to patient and physician preference. Treatment initiation is guided by recommendations issued by The American Association for the Study of Liver Diseases (AASLD) and the European Association for the Study of the Liver (EASL) and is based on detectable viral levels, HBeAg positive or negative status, ALT levels, and in certain cases, family history of HCC and liver biopsy. In patients with compensated cirrhosis, treatment is recommended regardless of HBeAg status or ALT level, but recommendations differ regarding HBV DNA levels; AASLD recommends treating at DNA levels detectable above 2x103 IU/mL; EASL and WHO recommend treating when HBV DNA levels are detectable at any level. In patients with decompensated cirrhosis, treatment and evaluation for liver transplantation are recommended in all cases if HBV DNA is detectable. Currently, multidrug treatment is not recommended in treatment of chronic HBV as it is no more effective in the long term than individual treatment with entecavir or tenofovir. Hepatitis C The American Association for the Study of Liver Diseases and the Infectious Diseases Society of America (AASLD-IDSA) recommend antiviral treatment for all patients with chronic hepatitis C infection except for those with additional chronic medical conditions that limit their life expectancy. Once it is acquired, persistence of the hepatitis C virus is the rule, resulting in chronic hepatitis C. The goal of treatment is prevention of hepatocellular carcinoma (HCC). The best way to reduce the long-term risk of HCC is to achieve sustained virological response (SVR). SVR is defined as an undetectable viral load at 12 weeks after treatment completion and indicates a cure. Currently available treatments include indirect and direct acting antiviral drugs. The indirect acting antivirals include pegylated interferon (PEG IFN) and ribavirin (RBV), which in combination have historically been the basis of therapy for HCV. Duration of and response to these treatments varies based on genotype. These agents are poorly tolerated but are still used in some resource-poor areas. In high-resource countries, they have been supplanted by direct acting antiviral agents, which first appeared in 2011; these agents target proteins responsible for viral replication and include the following three classes: NS3 and NS4A protease inhibitors, including telaprevir, boceprevir, simeprevir, and others NS5A inhibitors, including ledipasvir, daclatasvir, and others NS5B inhibitors, including sofosbuvir, dasabuvir, and others These drugs are used in various combinations, sometimes combined with ribavirin, based on the patient's genotype, delineated as genotypes 1–6. Genotype 1 (GT1), which is the most prevalent genotype in the United States and around the world, can now be cured with a direct acting antiviral regimen. First-line therapy for GT1 is a combination of sofosbuvir and ledipasvir (SOF/LDV) for 12 weeks for most patients, including those with advanced fibrosis or cirrhosis. Certain patients with early disease need only 8 weeks of treatment while those with advanced fibrosis or cirrhosis who have not responded to prior treatment require 24 weeks. Cost remains a major factor limiting access to these drugs, particularly in low-resource nations; the cost of the 12-week GT1 regimen (SOF/LDV) has been estimated at US$94,500. Hepatitis D Hepatitis D is difficult to treat, and effective treatments are lacking. Interferon alpha has proven effective at inhibiting viral activity but only on a temporary basis. Hepatitis E Similar to hepatitis A, treatment of hepatitis E is supportive and includes rest and ensuring adequate nutrition and hydration. Hospitalization may be required for particularly severe cases or for pregnant women. Alcoholic hepatitis First-line treatment of alcoholic hepatitis is treatment of alcoholism. For those who abstain completely from alcohol, reversal of liver disease and a longer life are possible; patients at every disease stage have been shown to benefit by prevention of additional liver injury. In addition to referral to psychotherapy and other treatment programs, treatment should include nutritional and psychosocial evaluation and treatment. Patients should also be treated appropriately for related signs and symptoms, such as ascites, hepatic encephalopathy, and infection. Severe alcoholic hepatitis has a poor prognosis and is notoriously difficult to treat. Without any treatment, 20–50% of patients may die within a month, but evidence shows treatment may extend life beyond one month (i.e., reduce short-term mortality). Available treatment options include pentoxifylline (PTX), which is a nonspecific TNF inhibitor, corticosteroids, such as prednisone or prednisolone (CS), corticosteroids with N-acetylcysteine (CS with NAC), and corticosteroids with pentoxifylline (CS with PTX). Data suggest that CS alone or CS with NAC are most effective at reducing short-term mortality. Unfortunately, corticosteroids are contraindicated in some patients, such as those who have active gastrointestinal bleeding, infection, kidney failure, or pancreatitis. In these cases PTX may be considered on a case-by-case basis in lieu of CS; some evidence shows PTX is better than no treatment at all and may be comparable to CS while other data show no evidence of benefit over placebo. Unfortunately, there are currently no drug treatments that decrease these patients' risk of dying in the longer term, at 3–12 months and beyond. Weak evidence suggests milk thistle extracts may improve survival in alcoholic liver disease and improve certain liver tests (serum bilirubin and GGT) without causing side effects, but a firm recommendation cannot be made for or against milk thistle without further study. The modified Maddrey's discriminant function may be used to evaluate the severity and prognosis in alcoholic hepatitis and evaluates the efficacy of using alcoholic hepatitis corticosteroid treatment. Metabolic hepatitis The main treatment of NASH is gradual weight loss and increased physical activity. In the United States, no medications have been approved to treat this disease. Autoimmune hepatitis Autoimmune hepatitis is commonly treated by immunosuppressants such as the corticosteroids prednisone or prednisolone, the active version of prednisolone that does not require liver synthesis, either alone or in combination with azathioprine, and some have suggested the combination therapy is preferred to allow for lower doses of corticosteroids to reduce associated side effects, although the result of treatment efficacy is comparative. Treatment of autoimmune hepatitis consists of two phases; an initial and maintenance phase. The initial phase consists of higher doses of corticosteroids which are tapered down over a number of weeks to a lower dose. If used in combination, azathioprine is given during the initial phase as well. Once the initial phase has been completed, a maintenance phase that consists of lower dose corticosteroids, and in combination therapy, azathioprine until liver blood markers are normalized. Treatment results in 66–91% of patients achieving normal liver test values in two years, with the average being 22 months. Prognosis Acute hepatitis Nearly all patients with hepatitis A infections recover completely without complications if they were healthy prior to the infection. Similarly, acute hepatitis B infections have a favorable course towards complete recovery in 95–99% of patients. Certain factors may portend a poorer outcome, such as co-morbid medical conditions or initial presenting symptoms of ascites, edema, or encephalopathy. Overall, the mortality rate for acute hepatitis is low: ~0.1% in total for cases of hepatitis A and B, but rates can be higher in certain populations (super infection with both hepatitis B and D, pregnant women, etc.). In contrast to hepatitis A & B, hepatitis C carries a much higher risk of progressing to chronic hepatitis, approaching 85–90%. Cirrhosis has been reported to develop in 20–50% of patients with chronic hepatitis C. Other rare complications of acute hepatitis include pancreatitis, aplastic anemia, peripheral neuropathy, and myocarditis. Fulminant hepatitis Despite the relatively benign course of most viral cases of hepatitis, fulminant hepatitis represents a rare but feared complication. Fulminant hepatitis most commonly occurs in hepatitis B, D, and E. About 1–2% of cases of hepatitis E can lead to fulminant hepatitis, but pregnant women are particularly susceptible, occurring in up to 20% of cases. Mortality rates in cases of fulminant hepatitis rise over 80%, but those patients that do survive often make a complete recovery. Liver transplantation can be life-saving in patients with fulminant liver failure. Hepatitis D infections can transform benign cases of hepatitis B into severe, progressive hepatitis, a phenomenon known as superinfection. Chronic hepatitis Acute hepatitis B infections become less likely to progress to chronic forms as the age of the patient increases, with rates of progression approaching 90% in vertically transmitted cases of infants compared to 1% risk in young adults. Overall, the five-year survival rate for chronic hepatitis B ranges from 97% in mild cases to 55% in severe cases with cirrhosis. Most patients who acquire hepatitis D at the same time as hepatitis B (co-infection) recover without developing a chronic infection. In people with hepatitis B who later acquire hepatitis D (superinfection), chronic infection is much more common at 80–90%, and liver disease progression is accelerated. Chronic hepatitis C progresses towards cirrhosis, with estimates of cirrhosis prevalence of 16% at 20 years after infection. While the major causes of mortality in hepatitis C is end stage liver disease, hepatocellular carcinoma is an important additional long term complication and cause of death in chronic hepatitis. Rates of mortality increase with progression of the underlying liver disease. Series of patients with compensated cirrhosis due to HCV have shown 3,5, and 10-year survival rates of 96, 91, and 79% respectively. The 5-year survival rate drops to 50% upon if the cirrhosis becomes decompensated. Epidemiology Viral hepatitis Hepatitis A Hepatitis A is found throughout the world and manifests as large outbreaks and epidemics associated with fecal contamination of water and food sources. Hepatitis A viral infection is predominant in children ages 5–14 with rare infection of infants. Infected children have little to no apparent clinical illness, in contrast to adults in whom greater than 80% are symptomatic if infected. Infection rates are highest in low resource countries with inadequate public sanitation and large concentrated populations. In such regions, as much as 90% of children younger than 10 years old have been infected and are immune, corresponding both to lower rates of clinically symptomatic disease and outbreaks. The availability of a childhood vaccine has significantly reduced infections in the United States, with incidence declining by more than 95% as of 2013. Paradoxically, the highest rates of new infection now occur in young adults and adults who present with worse clinical illness. Specific populations at greatest risk include: travelers to endemic regions, men who have sex with men, those with occupational exposure to non-human primates, people with clotting disorders who have received clotting factors, people with history of chronic liver disease in whom co-infection with hepatitis A can lead to fulminant hepatitis, and intravenous drug users (rare). Hepatitis B Hepatitis B is the most common cause of viral hepatitis in the world with more than 240 million chronic carriers of the virus, 1 million of whom are in the United States. In approximately two-thirds of patients who develop acute hepatitis B infection, no identifiable exposure is evident. Of those acutely infected, 25% become lifetime carriers of the virus. Risk of infection is highest among intravenous drug users, people with high-risk sexual behaviors, healthcare workers, people who had multiple transfusions, organ transplant patients, dialysis patients and newborns infected during the birthing process. Close to 780,000 deaths in the world are attributed to hepatitis B. The most endemic regions are in sub-Saharan Africa and East Asia, where as many as 10% of adults are chronic carriers. Carrier rates in developed nations are significantly lower, encompassing less than 1% of the population. In endemic regions, transmission is thought to be associated with exposure during birth and close contact between young infants. Hepatitis C Chronic hepatitis C is a major cause of liver cirrhosis and hepatocellular carcinoma. It is a common medical reason for liver transplantation due to its severe complications. It is estimated that 130–180 million people in the world are affected by this disease representing a little more than 3% of the world population. In the developing regions of Africa, Asia and South America, prevalence can be as high as 10% of the population. In Egypt, rates of hepatitis C infection as high as 20% have been documented and are associated with iatrogenic contamination related to schistosomiasis treatment in the 1950s–1980s. Currently in the United States, approximately 3.5 million adults are estimated to be infected. Hepatitis C is particularly prevalent among people born between 1945 and 1965, a group of about 800,000 people, with prevalence as high as 3.2% versus 1.6% in the general U.S. population. Most chronic carriers of hepatitis C are unaware of their infection status. The most common mode of transmission of hepatitis C virus is exposure to blood products via blood transfusions (prior to 1992) and intravenous drug injection. A history of intravenous drug injection is the most important risk factor for chronic hepatitis C. Other susceptible populations include those engaged in high-risk sexual behavior, infants of infected mothers, and healthcare workers. Hepatitis D The hepatitis D virus causes chronic and fulminant hepatitis in the context of co-infection with the hepatitis B virus. It is primarily transmitted via non-sexual contact and via needles. Susceptibility to hepatitis D differs by geographic region. In the United States and Northern Europe, populations at risk are intravenous drug users and people who receive multiple transfusions. In the Mediterranean, hepatitis D is predominant among hepatitis B virus co-infected people. Hepatitis E Similar to Hepatitis A, hepatitis E manifests as large outbreaks and epidemics associated with fecal contamination of water sources. It accounts for more than 55,000 deaths annually with approximately 20 million people worldwide thought to be infected with the virus. It affects predominantly young adults, causing acute hepatitis. In infected pregnant women, Hepatitis E infection can lead to fulminant hepatitis with third trimester mortality rates as high as 30%. Those with weakened immune systems, such as organ transplant recipients, are also susceptible. Infection is rare in the United States but rates are high in the developing world (Africa, Asia, Central America, Middle East). Many genotypes exist and are differentially distributed around the world. There is some evidence of hepatitis E infection of animals, serving as a reservoir for human infection. Alcoholic hepatitis Alcoholic hepatitis (AH) in its severe form has a one-month mortality as high as 50%. Most people who develop AH are men but women are at higher risk of developing AH and its complications likely secondary to high body fat and differences in alcohol metabolism. Other contributing factors include younger age <60, binge pattern drinking, poor nutritional status, obesity and hepatitis C co-infection. It is estimated that as much as 20% of people with AH are also infected with hepatitis C. In this population, the presence of hepatitis C virus leads to more severe disease with faster progression to cirrhosis, hepatocellular carcinoma and increased mortality. Obesity increases the likelihood of progression to cirrhosis in cases of alcoholic hepatitis. It is estimated that 70% of people who have AH will progress to cirrhosis. Non-alcoholic steatohepatitis Non-alcoholic steatohepatitis (NASH) is projected to become the top reason for liver transplantation in the United States by 2020, supplanting chronic liver disease due to hepatitis C. About 20–45% of the U.S. population have NAFLD and 6% have NASH. The estimated prevalence of NASH in the world is 3–5%. Of NASH patients who develop cirrhosis, about 2% per year will likely progress to hepatocellular carcinoma. Worldwide, the estimated prevalence of hepatocellular carcinoma related to NAFLD is 15–30%. NASH is thought to be the primary cause of cirrhosis in approximately 25% of patients in the United States, representing 1–2% of the general population. History Early observations Initial accounts of a syndrome that we now think is likely to be hepatitis begin to occur around 3000 B.C. Clay tablets that served as medical handbooks for the ancient Sumerians described the first observations of jaundice. The Sumerians believed that the liver was the home of the soul, and attributed the findings of jaundice to the attack of the liver by a devil named Ahhazu. Around 400 B.C., Hippocrates recorded the first documentation of an epidemic jaundice, in particular noting the uniquely fulminant course of a cohort of patients who all died within two weeks. He wrote, "The bile contained in the liver is full of phlegm and blood, and erupts...After such an eruption, the patient soon raves, becomes angry, talks nonsense and barks like a dog." Given the poor sanitary conditions of war, infectious jaundice played a large role as a major cause of mortality among troops in the Napoleonic Wars, the American Revolutionary War, and both World Wars. During World War II, estimates of soldiers affected by hepatitis were upwards of 10 million. During World War II, soldiers received vaccines against diseases such as yellow fever, but these vaccines were stabilized with human serum, presumably contaminated with hepatitis viruses, which often created epidemics of hepatitis. It was suspected these epidemics were due to a separate infectious agent, and not due to the yellow fever virus itself, after noting 89 cases of jaundice in the months after vaccination out of a total 3,100 patients that were vaccinated. After changing the seed virus strain, no cases of jaundice were observed in the subsequent 8,000 vaccinations. Willowbrook State School experiments A New York University researcher named Saul Krugman continued this research into the 1950s and 1960s, most infamously with his experiments on mentally disabled children at the Willowbrook State School in New York, a crowded urban facility where hepatitis infections were highly endemic to the student body. Krugman injected students with gamma globulin, a type of antibody. After observing the temporary protection against infection this antibody provided, he then tried injected live hepatitis virus into students. Krugman also controversially took feces from infected students, blended it into milkshakes, and fed it to newly admitted children. His research was received with much controversy, as people protested the questionable ethics surrounding the chosen target population. Henry Beecher was one of the foremost critics in an article in the New England Journal of Medicine in 1966, arguing that parents were unaware to the risks of consent and that the research was done to benefit others at the expense of children. Moreover, he argued that poor families with mentally disabled children often felt pressured to join the research project to gain admission to the school, with all of the educational and support resources that would come along with it. Others in the medical community spoke out in support of Krugman's research in terms of its widespread benefits and understanding of the hepatitis virus, and Willowbrook continues to be a commonly cited example in debates about medical ethics. Australia antigen The next insight regarding hepatitis B was a serendipitous one by Dr. Baruch Blumberg, a researcher at the NIH who did not set out to research hepatitis, but rather studied lipoprotein genetics. He travelled across the globe collecting blood samples, investigating the interplay between disease, environment, and genetics with the goal of designing targeted interventions for at-risk people that could prevent them from getting sick. He noticed an unexpected interaction between the blood of a patient with hemophilia that had received multiple transfusions and a protein found in the blood of an indigenous Australian person. He named the protein the "Australia antigen" and made it the focus of his research. He found a higher prevalence of the protein in the blood of patients from developing countries, compared to those from developed ones, and noted associations of the antigen with other diseases like leukemia and Down Syndrome. Eventually, he came to the unifying conclusion that the Australia antigen was associated with viral hepatitis. In 1970, David Dane first isolated the hepatitis B virion at London's Middlesex Hospital, and named the virion the 42-nm "Dane particle". Based on its association with the surface of the hepatitis B virus, the Australia antigen was renamed to "hepatitis B surface antigen" or HBsAg. Blumberg continued to study the antigen, and eventually developed the first hepatitis B vaccine using plasma rich in HBsAg, for which he received the Nobel Prize in Medicine in 1976. Society and culture Economic burden Overall, hepatitis accounts for a significant portion of healthcare expenditures in both developing and developed nations, and is expected to rise in several developing countries. While hepatitis A infections are self-limited events, they are associated with significant costs in the United States. It has been estimated that direct and indirect costs are approximately $1817 and $2459 respectively per case, and that an average of 27 work days is lost per infected adult. A 1997 report demonstrated that a single hospitalization related to hepatitis A cost an average of $6,900 and resulted in around $500 million in total annual healthcare costs. Cost effectiveness studies have found widespread vaccination of adults to not be feasible, but have stated that a combination hepatitis A and B vaccination of children and at risk groups (people from endemic areas, healthcare workers) may be. Hepatitis B accounts for a much larger percentage of health care spending in endemic regions like Asia. In 1997 it accounted for 3.2% of South Korea's total health care expenditures and resulted in $696 million in direct costs. A large majority of that sum was spent on treating disease symptoms and complications. Chronic hepatitis B infections are not as endemic in the United States, but accounted for $357 million in hospitalization costs in the year 1990. That number grew to $1.5 billion in 2003, but remained stable as of 2006, which may be attributable to the introduction of effective drug therapies and vaccination campaigns. People infected with chronic hepatitis C tend to be frequent users of the health care system globally. It has been estimated that a person infected with hepatitis C in the United States will result in a monthly cost of $691. That number nearly doubles to $1,227 for people with compensated (stable) cirrhosis, while the monthly cost of people with decompensated (worsening) cirrhosis is almost five times as large at $3,682. The wide-ranging effects of hepatitis make it difficult to estimate indirect costs, but studies have speculated that the total cost is $6.5 billion annually in the United States. In Canada, 56% of HCV related costs are attributable to cirrhosis and total expenditures related to the virus are expected to peak at CAD$396 million in the year 2032. 2003 Monaca outbreak The largest outbreak of hepatitis A virus in United States history occurred among people who ate at a now-defunct Mexican food restaurant located in Monaca, Pennsylvania in late 2003. Over 550 people who visited the restaurant between September and October 2003 were infected with the virus, three of whom died as a direct result. The outbreak was brought to the attention of health officials when local emergency medicine physicians noticed a significant increase in cases of hepatitis A in the county. After conducting its investigation, the CDC attributed the source of the outbreak to the use of contaminated raw green onion. The restaurant was purchasing its green onion stock from farms in Mexico at the time. It is believed that the green onions may have been contaminated through the use of contaminated water for crop irrigation, rinsing, or icing or by handling of the vegetables by infected people. Green onion had caused similar outbreaks of hepatitis A in the southern United States prior to this, but not to the same magnitude. The CDC believes that the restaurant's use of a large communal bucket for chopped raw green onion allowed non-contaminated plants to be mixed with contaminated ones, increasing the number of vectors of infection and amplifying the outbreak. The restaurant was closed once it was discovered to be the source, and over 9,000 people were given hepatitis A immune globulin because they had either eaten at the restaurant or had been in close contact with someone who had. Special populations HIV co-infection Persons infected with HIV have a particularly high burden of HIV-HCV co-infection. In a recent study by the WHO, the likelihood of being infected with hepatitis C virus was six times greater in those who also had HIV. The prevalence of HIV-HCV co-infection worldwide was estimated to be 6.2% representing more than 2.2 million people. Intravenous drug use was an independent risk factor for HCV infection. In the WHO study, the prevalence of HIV-HCV co-infection was markedly higher at 82.4% in those who injected drugs compared to the general population (2.4%). In a study of HIV-HCV co-infection among HIV positive men who have sex with men (MSM), the overall prevalence of anti-hepatitis C antibodies was estimated to be 8.1% and increased to 40% among HIV positive MSM who also injected drugs. Pregnancy Hepatitis B Vertical transmission is a significant contributor of new HBV cases each year, with 35–50% of transmission from mother to neonate in endemic countries. Vertical transmission occurs largely via a neonate's exposure to maternal blood and vaginal secretions during birth. While the risk of progression to chronic infection is approximately 5% among adults who contract the virus, it is as high as 95% among neonates subject to vertical transmission. The risk of viral transmission is approximately 10–20% when maternal blood is positive for HBsAg, and up to 90% when also positive for HBeAg. Given the high risk of perinatal transmission, the CDC recommends screening all pregnant women for HBV at their first prenatal visit. It is safe for non-immune pregnant women to receive the HBV vaccine. Based on the limited available evidence, the American Association for the Study of Liver Diseases (AASLD) recommends antiviral therapy in pregnant women whose viral load exceeds 200,000 IU/mL. A growing body of evidence shows that antiviral therapy initiated in the third trimester significantly reduces transmission to the neonate. A systematic review of the Antiretroviral Pregnancy Registry database found that there was no increased risk of congenital anomalies with Tenofovir; for this reason, along with its potency and low risk of resistance, the AASLD recommends this drug. A 2010 systematic review and meta-analysis found that Lamivudine initiated early in the third trimester also significantly reduced mother-to-child transmission of HBV, without any known adverse effects. The ACOG states that the evidence available does not suggest any particular mode of delivery (i.e. vaginal vs. cesarean) is better at reducing vertical transmission in mothers with HBV. The WHO and CDC recommend that neonates born to mothers with HBV should receive hepatitis B immune globulin (HBIG) as well as the HBV vaccine within 12 hours of birth. For infants who have received HBIG and the HBV vaccine, breastfeeding is safe. Hepatitis C Estimates of the rate of HCV vertical transmission range from 2–8%; a 2014 systematic review and meta-analysis found the risk to be 5.8% in HCV-positive, HIV-negative women. The same study found the risk of vertical transmission to be 10.8% in HCV-positive, HIV-positive women. Other studies have found the risk of vertical transmission to be as high as 44% among HIV-positive women. The risk of vertical transmission is higher when the virus is detectable in the mother's blood. Evidence does not indicate that mode of delivery (i.e. vaginal vs. cesarean) has an effect on vertical transmission. For women who are HCV-positive and HIV-negative, breastfeeding is safe. CDC guidelines suggest avoiding it if a woman's nipples are cracked or bleeding to reduce the risk of transmission. Hepatitis E Pregnant women who contract HEV are at significant risk of developing fulminant hepatitis with maternal mortality rates as high as 20–30%, most commonly in the third trimester . A 2016 systematic review and meta-analysis of 47 studies that included 3968 people found maternal case-fatality rates (CFR) of 20.8% and fetal CFR of 34.2%; among women who developed fulminant hepatic failure, CFR was 61.2%. Vaccine An essential defense against hepatitis infections, especially those caused by Hepatitis A and B, is the hepatitis vaccination. The Hepatitis B vaccination is quite effective and frequently used. The frequency of hepatitis-related liver illnesses and fatalities has been considerably decreased by the immunization campaigns.	Biology and health sciences	Non-infectious disease	null
38257	https://en.wikipedia.org/wiki/Electrolysis	Electrolysis	In chemistry and manufacturing, electrolysis is a technique that uses direct electric current (DC) to drive an otherwise non-spontaneous chemical reaction. Electrolysis is commercially important as a stage in the separation of elements from naturally occurring sources such as ores using an electrolytic cell. The voltage that is needed for electrolysis to occur is called the decomposition potential. The word "lysis" means to separate or break, so in terms, electrolysis would mean "breakdown via electricity." Etymology The word "electrolysis" was introduced by Michael Faraday in 1834, using the Greek words "amber", which since the 17th century was associated with electrical phenomena, and meaning "dissolution". Nevertheless, electrolysis, as a tool to study chemical reactions and obtain pure elements, precedes the coinage of the term and formal description by Faraday. History In the early nineteenth century, William Nicholson and Anthony Carlisle sought to further Volta's experiments. They attached two wires to either side of a voltaic pile and placed the other ends in a tube filled with water. They noticed when the wires were brought together that each wire produced bubbles. One type was hydrogen, the other was oxygen. In 1785 a Dutch scientist named Martin van Marum created an electrostatic generator that he used to reduce tin, zinc and antimony from their salts using a process later known as electrolysis. Though he unknowingly produced electrolysis, it was not until 1800 when William Nicholson and Anthony Carlisle discovered how electrolysis works. In 1791 Luigi Galvani experimented with frog legs. He claimed that placing animal muscle between two dissimilar metal sheets resulted in electricity. Responding to these claims, Alessandro Volta conducted his own tests. This would give insight to Humphry Davy's ideas on electrolysis. During preliminary experiments, Humphry Davy hypothesized that when two elements combine to form a compound, electrical energy is released. Humphry Davy would go on to create Decomposition Tables from his preliminary experiments on Electrolysis. The Decomposition Tables would give insight on the energies needed to break apart certain compounds. In 1817 Johan August Arfwedson determined there was another element, lithium, in some of his samples; however, he could not isolate the component. It was not until 1821 that William Thomas Brande used electrolysis to single it out. Two years later, he streamlined the process using lithium chloride and potassium chloride with electrolysis to produce lithium and lithium hydroxide. During the later years of Humphry Davy's research, Michael Faraday became his assistant. While studying the process of electrolysis under Humphry Davy, Michael Faraday discovered two laws of electrolysis. During the time of Maxwell and Faraday, concerns came about for electropositive and electronegative activities. In November 1875, Paul Émile Lecoq de Boisbaudran discovered gallium using electrolysis of gallium hydroxide, producing 3.4 mg of gallium. The following December, he presented his discovery of gallium to the Académie des sciences in Paris. On June 26, 1886, Ferdinand Frederick Henri Moissan finally felt comfortable performing electrolysis on anhydrous hydrogen fluoride to create a gaseous fluorine pure element. Before he used hydrogen fluoride, Henri Moissan used fluoride salts with electrolysis. Thus on June 28, 1886, he performed his experiment in front of the Académie des sciences to show his discovery of the new element fluorine. While trying to find elemental fluorine through electrolysis of fluoride salts, many chemists perished including Paulin Louyet and Jérôme Nicklès. In 1886 Charles Martin Hall from America and Paul Héroult from France both filed for American patents for the electrolysis of aluminum, with Héroult submitting his in May, and Hall, in July. Hall was able to get his patent by proving through letters to his brother and family evidence that his method was discovered before the French patent was submitted. This became known as the Hall–Héroult process which benefited many industries because aluminum's price then dropped from four dollars to thirty cents per pound. In 1902 Polish engineer and inventor Stanisław Łaszczyński filed for and obtained Polish patent for the electrolysis of copper and zinc. Timeline 1785 – Martinus van Marum's electrostatic generator was used to reduce tin, zinc, and antimony from their salts using electrolysis. 1800 – William Nicholson and Anthony Carlisle (and also Johann Ritter), decomposed water into hydrogen and oxygen. 1808 – Potassium (1807), sodium (1807), barium, calcium and magnesium were discovered by Humphry Davy using electrolysis. 1821 – Lithium was discovered by the English chemist William Thomas Brande, who obtained it by electrolysis of lithium oxide. 1834 – Michael Faraday published his two laws of electrolysis, provided a mathematical explanation for them, and introduced terminology such as electrode, electrolyte, anode, cathode, anion, and cation. 1875 – Paul Émile Lecoq de Boisbaudran discovered gallium using electrolysis. 1886 – Fluorine was discovered by Henri Moissan using electrolysis. 1886 – Hall–Héroult process developed for making aluminium. 1890 – Castner–Kellner process developed for making sodium hydroxide. 1902 – Stanisław Łaszczyński obtained copper using electrolysis. Overview Electrolysis is the passing of a direct electric current through an electrolyte which is producing chemical reactions at the electrodes and decomposition of the materials. The main components required to achieve electrolysis are an electrolyte, electrodes, and an external power source. A partition (e.g. an ion-exchange membrane or a salt bridge) is optional to keep the products from diffusing to the vicinity of the opposite electrode. The electrolyte is a chemical substance which contains free ions and carries electric current (e.g. an ion-conducting polymer, solution, or a ionic liquid compound). If the ions are not mobile, as in most solid salts, then electrolysis cannot occur. A liquid electrolyte is produced by: Solvation or reaction of an ionic compound with a solvent (such as water) to produce mobile ions An ionic compound melted by heating The electrodes are immersed separated by a distance such that a current flows between them through the electrolyte and are connected to the power source which completes the electrical circuit. A direct current supplied by the power source drives the reaction causing ions in the electrolyte to be attracted toward the respective oppositely charged electrode. Electrodes of metal, graphite and semiconductor material are widely used. Choice of suitable electrode depends on chemical reactivity between the electrode and electrolyte and manufacturing cost. Historically, when non-reactive anodes were desired for electrolysis, graphite (called plumbago in Faraday's time) or platinum were chosen. They were found to be some of the least reactive materials for anodes. Platinum erodes very slowly compared to other materials, and graphite crumbles and can produce carbon dioxide in aqueous solutions but otherwise does not participate in the reaction. Cathodes may be made of the same material, or they may be made from a more reactive one since anode wear is greater due to oxidation at the anode. Process of electrolysis The key process of electrolysis is the interchange of atoms and ions by the removal or addition of electrons due to the applied potential. The desired products of electrolysis are often in a different physical state from the electrolyte and can be removed by mechanical processes (e.g. by collecting gas above an electrode or precipitating a product out of the electrolyte). The quantity of the products is proportional to the current, and when two or more electrolytic cells are connected in series to the same power source, the products produced in the cells are proportional to their equivalent weight. These are known as Faraday's laws of electrolysis. Each electrode attracts ions that are of the opposite charge. Positively charged ions (cations) move towards the electron-providing (negative) cathode. Negatively charged ions (anions) move towards the electron-extracting (positive) anode. In this process electrons are effectively introduced at the cathode as a reactant and removed at the anode as a product. In chemistry, the loss of electrons is called oxidation, while electron gain is called reduction. When neutral atoms or molecules, such as those on the surface of an electrode, gain or lose electrons they become ions and may dissolve in the electrolyte and react with other ions. When ions gain or lose electrons and become neutral, they will form compounds that separate from the electrolyte. Positive metal ions like Cu2+ deposit onto the cathode in a layer. The terms for this are electroplating, electrowinning, and electrorefining. When an ion gains or loses electrons without becoming neutral, its electronic charge is altered in the process. For example, the electrolysis of brine produces hydrogen and chlorine gases which bubble from the electrolyte and are collected. The initial overall reaction is thus: 2 NaCl + 2 H2O → 2 NaOH + H2 + Cl2 The reaction at the anode results in chlorine gas from chlorine ions: 2 Cl− → Cl2 + 2 e− The reaction at the cathode results in hydrogen gas and hydroxide ions: 2 H2O + 2 e− → H2 + 2 OH− Without a partition between the electrodes, the OH− ions produced at the cathode are free to diffuse throughout the electrolyte to the anode. As the electrolyte becomes more basic due to the production of OH−, less Cl2 emerges from the solution as it begins to react with the hydroxide producing hypochlorite (ClO−) at the anode: Cl2 + 2 NaOH → NaCl + NaClO + H2O The more opportunity the Cl2 has to interact with NaOH in the solution, the less Cl2 emerges at the surface of the solution and the faster the production of hypochlorite progresses. This depends on factors such as solution temperature, the amount of time the Cl2 molecule is in contact with the solution, and concentration of NaOH. Likewise, as hypochlorite increases in concentration, chlorates are produced from them: 3 NaClO → NaClO3 + 2 NaCl Other reactions occur, such as the self-ionization of water and the decomposition of hypochlorite at the cathode, the rate of the latter depends on factors such as diffusion and the surface area of the cathode in contact with the electrolyte. Decomposition potential Decomposition potential or decomposition voltage refers to the minimum voltage (difference in electrode potential) between anode and cathode of an electrolytic cell that is needed for electrolysis to occur. The voltage at which electrolysis is thermodynamically preferred is the difference of the electrode potentials as calculated using the Nernst equation. Applying additional voltage, referred to as overpotential, can increase the rate of reaction and is often needed above the thermodynamic value. It is especially necessary for electrolysis reactions involving gases, such as oxygen, hydrogen or chlorine. Oxidation and reduction at the electrodes Oxidation of ions or neutral molecules occurs at the anode. For example, it is possible to oxidize ferrous ions to ferric ions at the anode: Fe(aq) → Fe(aq) + e− Reduction of ions or neutral molecules occurs at the cathode. It is possible to reduce ferricyanide ions to ferrocyanide ions at the cathode: Fe(CN) + e− → Fe(CN) Neutral molecules can also react at either of the electrodes. For example: p-benzoquinone can be reduced to hydroquinone at the cathode: + 2 e− + 2 H+ → In the last example, H+ ions (hydrogen ions) also take part in the reaction and are provided by the acid in the solution, or by the solvent itself (water, methanol, etc.). Electrolysis reactions involving H+ ions are fairly common in acidic solutions. In aqueous alkaline solutions, reactions involving OH− (hydroxide ions) are common. Sometimes the solvents themselves (usually water) are oxidized or reduced at the electrodes. It is even possible to have electrolysis involving gases, e.g. by using a gas diffusion electrode. Energy changes during electrolysis The amount of electrical energy that must be added equals the change in Gibbs free energy of the reaction plus the losses in the system. The losses can (in theory) be arbitrarily close to zero, so the maximum thermodynamic efficiency equals the enthalpy change divided by the free energy change of the reaction. In most cases, the electric input is larger than the enthalpy change of the reaction, so some energy is released in the form of heat. In some cases, for instance, in the electrolysis of steam into hydrogen and oxygen at high temperature, the opposite is true and heat energy is absorbed. This heat is absorbed from the surroundings, and the heating value of the produced hydrogen is higher than the electric input. Variations Pulsating current results in products different from DC. For example, pulsing increases the ratio of ozone to oxygen produced at the anode in the electrolysis of an aqueous acidic solution such as dilute sulphuric acid. Electrolysis of ethanol with pulsed current evolves an aldehyde instead of primarily an acid. Related processes Galvanic cells and batteries use spontaneous, energy-releasing redox reactions to generate an electrical potential that provides useful power. When a secondary battery is charged, its redox reaction is run in reverse and the system can be considered as an electrolytic cell. Industrial uses Chloralkali process The chloralkali process is a large scale application of electrolysis. This technology supplies most of the chlorine and sodium hydroxide required by many industries. The cathode is a mixed metal oxide clad titanium anode (also called a dimensionally stable anode). Electrofluorination Many organofluorine compounds are produced by electrofluorination. One manifestation of this technology is the Simons process, which can be described as: R3C–H + HF → R3C–F + H2 In the course of a typical synthesis, this reaction occurs once for each C–H bond in the precursor. The cell potential is maintained near 5–6 V. The anode, the electrocatalyst, is nickel-plated. Hydrodimerization of acrylonitrile Acrylonitrile is converted to adiponitrile on an industrial scale via electrocatalysis. Electroplating and electrowinning processes Purifying copper from refined copper. Electrometallurgy of aluminium, lithium, sodium, potassium, magnesium, calcium. Electroplating, where a thin film of metal is deposited over a substrate material. Electroplating is used in many industries for either functional or decorative purposes, as in-vehicle bodies and nickel coins. Electrochemical machining (ECM) In Electrochemical machining, an electrolytic cathode is used as a shaped tool for removing material by anodic oxidation from a workpiece. ECM is often used as a technique for deburring or for etching metal surfaces like tools or knives with a permanent mark or logo. Other Production of sodium chlorate and potassium chlorate. Production of fuels such as hydrogen for spacecraft, nuclear submarines and vehicles. Rust removal and cleaning of old coins and other metallic objects. Competing half-reactions in solution electrolysis Using a cell containing inert platinum electrodes, electrolysis of aqueous solutions of some salts leads to the reduction of the cations (such as metal deposition with, for example, zinc salts) and oxidation of the anions (such as the evolution of bromine with bromides). However, with salts of some metals (such as sodium) hydrogen is evolved at the cathode, and for salts containing some anions (such as sulfate ) oxygen is evolved at the anode. In both cases, this is due to water being reduced to form hydrogen or oxidized to form oxygen. In principle, the voltage required to electrolyze a salt solution can be derived from the standard electrode potential for the reactions at the anode and cathode. The standard electrode potential is directly related to the Gibbs free energy, ΔG, for the reactions at each electrode and refers to an electrode with no current flowing. An extract from the table of standard electrode potentials is shown below. {\| class="wikitable" \|- ! Half-reaction ! E° (V) ! Ref. \|- \| Na+ + e− Na \| −2.71 \|\| \|- \| Zn2+ + 2 e− Zn \| −0.7618 \|\| \|- \| 2 H+ + 2 e− H2 \| ≡ 0\|\| \|- \| Br2 + 2 e− 2 Br− \| +1.0873 \|\| \|- \| O2 + 4 H+ + 4 e− 2 H2O \| +1.23 \|\| \|- \| Cl2 + 2 e− 2 Cl− \| +1.36 \|\| \|- \| + 2 e− 2 \| +2.07 \|\| \|} In terms of electrolysis, this table should be interpreted as follows: Moving down the table, E° becomes more positive, and species on the left are more likely to be reduced: for example, zinc ions are more likely to be reduced to zinc metal than sodium ions are to be reduced to sodium metal. Moving up the table, E° becomes more negative, and species on the right are more likely to be oxidized: for example, sodium metal is more likely to be oxidized to sodium ions than zinc metal is to be oxidized to zinc ions. Using the Nernst equation the electrode potential can be calculated for a specific concentration of ions, temperature and the number of electrons involved. For pure water (pH 7): the electrode potential for the reduction producing hydrogen is −0.41 V, the electrode potential for the oxidation producing oxygen is +0.82 V. Comparable figures calculated in a similar way, for 1 M zinc bromide, ZnBr2, are −0.76 V for the reduction to Zn metal and +1.10 V for the oxidation producing bromine. The conclusion from these figures is that hydrogen should be produced at the cathode and oxygen at the anode from the electrolysis of water—which is at variance with the experimental observation that zinc metal is deposited and bromine is produced. The explanation is that these calculated potentials only indicate the thermodynamically preferred reaction. In practice, many other factors have to be taken into account such as the kinetics of some of the reaction steps involved. These factors together mean that a higher potential is required for the reduction and oxidation of water than predicted, and these are termed overpotentials. Experimentally it is known that overpotentials depend on the design of the cell and the nature of the electrodes. For the electrolysis of a neutral (pH 7) sodium chloride solution, the reduction of sodium ion is thermodynamically very difficult and water is reduced evolving hydrogen leaving hydroxide ions in solution. At the anode the oxidation of chlorine is observed rather than the oxidation of water since the overpotential for the oxidation of chloride to chlorine is lower than the overpotential for the oxidation of water to oxygen. The hydroxide ions and dissolved chlorine gas react further to form hypochlorous acid. The aqueous solutions resulting from this process is called electrolyzed water and is used as a disinfectant and cleaning agent. Research trends Electrolysis of carbon dioxide The electrochemical reduction or electrocatalytic conversion of CO2 can produce value-added chemicals such as methane, ethylene, ethanol, etc. The electrolysis of carbon dioxide gives formate or carbon monoxide, but sometimes more elaborate organic compounds such as ethylene. This technology is under research as a carbon-neutral route to organic compounds. Electrolysis of acidified water Electrolysis of water produces hydrogen and oxygen in a ratio of 2 to 1 respectively. 2 H2O → 2 H2 + O2 E° = +1.229 V The energy efficiency of water electrolysis varies widely. The efficiency of an electrolyser is a measure of the enthalpy contained in the hydrogen (to undergo combustion with oxygen or some other later reaction), compared with the input electrical energy. Heat/enthalpy values for hydrogen are well published in science and engineering texts, as 144 MJ/kg (40 kWh/kg). Note that fuel cells (not electrolysers) cannot use this full amount of heat/enthalpy, which has led to some confusion when calculating efficiency values for both types of technology. In the reaction, some energy is lost as heat. Some reports quote efficiencies between 50% and 70% for alkaline electrolysers (50 kWh/kg); however, higher practical efficiencies are available with the use of polymer electrolyte membrane electrolysis and catalytic technology, such as 95% efficiency. The National Renewable Energy Laboratory estimated in 2006 that 1 kg of hydrogen (roughly equivalent to 3 kg, or 4 liters, of petroleum in energy terms) could be produced by wind powered electrolysis for between US$5.55 in the near term and US$2.27 in the longer term. About 4% of hydrogen gas produced worldwide is generated by electrolysis, and normally used onsite. Hydrogen is used for the creation of ammonia for fertilizer via the Haber process, and converting heavy petroleum sources to lighter fractions via hydrocracking. Onsite electrolysis has been utilized to capture hydrogen for hydrogen fuel-cells in hydrogen vehicles. Carbon/hydrocarbon assisted water electrolysis Recently, to reduce the energy input, the utilization of carbon (coal), alcohols (hydrocarbon solution), and organic solution (glycerol, formic acid, ethylene glycol, etc.) with co-electrolysis of water has been proposed as a viable option. The carbon/hydrocarbon assisted water electrolysis (so-called CAWE) process for hydrogen generation would perform this operation in a single electrochemical reactor. This system energy balance can be required only around 40% electric input with 60% coming from the chemical energy of carbon or hydrocarbon. This process utilizes solid coal/carbon particles or powder as fuels dispersed in acid/alkaline electrolyte in the form of slurry and the carbon contained source co-assist in the electrolysis process as following theoretical overall reactions: Carbon/Coal slurry (C + 2H2O) → CO2 + 2H2 E′ = 0.21 V (reversible voltage) / E′ = 0.46 V (thermo-neutral voltage) or Carbon/Coal slurry (C + H2O) → CO + H2 E′ = 0.52 V (reversible voltage) / E′ = 0.91 V (thermo-neutral voltage) Thus, this CAWE approach is that the actual cell overpotential can be significantly reduced to below 1.0 V as compared to 1.5 V for conventional water electrolysis. Electrocrystallization A specialized application of electrolysis involves the growth of conductive crystals on one of the electrodes from oxidized or reduced species that are generated in situ. The technique has been used to obtain single crystals of low-dimensional electrical conductors, such as charge-transfer salts and linear chain compounds. Electrolysis of Iron Ore The current method of producing steel from iron ore is very carbon intensive, in part to the direct release of CO2 in the blast furnace. A study of steel making in Germany found that producing 1 ton of steel emitted 2.1 tons of CO2e with 22% of that being direct emissions from the blast furnace. As of 2022, steel production contributes 7–9% of global emissions. Electrolysis of iron can eliminate direct emissions and further reduce emissions if the electricity is created from green energy. The small-scale electrolysis of iron has been successfully reported by dissolving it in molten oxide salts and using a platinum anode. Oxygen anions form oxygen gas and electrons at the anode. Iron cations consume electrons and form iron metal at the cathode. This method was performed a temperature of 1550 °C which presents a significant challenge to maintaining the reaction. Particularly, anode corrosion is a concern at these temperatures. Additionally, the low temperature reduction of iron oxide by dissolving it in alkaline water has been reported. The temperature is much lower than traditional iron production at 114 °C. The low temperatures also tend to correlate with higher current efficiencies, with an efficiency of 95% being reported. While these methods are promising, they struggle to be cost competitive because of the large economies of scale keeping the price of blast furnace iron low. Electrolysis of seawater A 2020 study investigated direct electrolysis of seawater, alkaline electrolysis, proton-exchange membrane electrolysis, and solid oxide electrolysis. Direct electrolysis of seawater follows known processes, forming an electrolysis cell in which the seawater acts as the electrolyte to allow for the reaction at the anode, and the reaction at the cathode, . The inclusion of magnesium and calcium ions in the seawater makes the production of alkali hydroxides possible that could form scales in the electrolyser cell, cutting down on lifespan and increasing the need for maintenance. The alkaline electrolysers operate with the following reactions at the anode, and cathode, , and use high base solutions as electrolytes, operating at and need additional separators to ensure the gas phase hydrogen and oxygen remain separate. The electrolyte can easily get contaminated, but the alkaline electrolyser can operate under pressure to improve energy consumption. The electrodes can be made of inexpensive materials and there's no requirement for an expensive catalyst in the design. Proton-exchange membrane electrolysers operate with the reactions at the anode, and cathode, , at temperatures of , using a solid polymer electrolyte and requiring higher costs of processing to allow the solid electrolyte to touch uniformly to the electrodes. Similar to the alkaline electrolyser, the proton exchange membrane electrolyser can operate at higher pressures, reducing the energy costs required to compress the hydrogen gas afterward, but the proton exchange membrane electrolyser also benefits from rapid response times to changes in power requirements or demands and not needing maintenance, at the cost of having a faster inherent degradation rate and being the most vulnerable to impurities in the water. Solid oxide electrolysers run the reactions at the anode and at the cathode.The solid oxide electrolysers require high temperatures () to operate, generating superheated steam. They suffer from degradation when turned off, making it a more inflexible hydrogen generation technology. In a selected series of multiple-criteria decision-analysis comparisons in which the highest priority was placed on economic operation costs followed equally by environmental and social criteria, it was found that the proton exchange membrane electrolyser offered the most suitable combination of values (e.g., investment cost, maintenance, and operation cost, resistance to impurities, specific energy for hydrogen production at sea, risk of environmental impact, etc.), followed by the alkaline electrolyser, with the alkaline electrolyser being the most economically feasible, but more hazardous in terms of safety and environmental concerns due to the need for basic electrolyte solutions as opposed to the solid polymers used in proton-exchange membranes. Due to the methods conducted in multiple-criteria decision analysis, non-objective weights are applied to the various factors, and so multiple methods of decision analysis were performed simultaneously to examine the electrolysers in a way that minimizes the effects of bias on the performance conclusions.	Physical sciences	Chemistry: General	null
38272	https://en.wikipedia.org/wiki/Nicotine	Nicotine	Nicotine is a naturally produced alkaloid in the nightshade family of plants (most predominantly in tobacco and Duboisia hopwoodii) and is widely used recreationally as a stimulant and anxiolytic. As a pharmaceutical drug, it is used for smoking cessation to relieve withdrawal symptoms. Nicotine acts as a receptor agonist at most nicotinic acetylcholine receptors (nAChRs), except at two nicotinic receptor subunits (nAChRα9 and nAChRα10) where it acts as a receptor antagonist. Nicotine constitutes approximately 0.6–3.0% of the dry weight of tobacco. Nicotine is also present at ppb concentrations in edible plants in the family Solanaceae, including potatoes, tomatoes, and eggplants, though sources disagree on whether this has any biological significance to human consumers. It functions as an antiherbivore toxin; consequently, nicotine was widely used as an insecticide in the past, and neonicotinoids (structurally similar to nicotine), such as imidacloprid, are some of the most effective and widely used insecticides. Nicotine is highly addictive. Slow-release forms (gums and patches, when used correctly) can be less addictive and help in quitting. Animal research suggests that monoamine oxidase inhibitors present in tobacco smoke may enhance nicotine's addictive properties. An average cigarette yields about 2 mg of absorbed nicotine. The estimated lower dose limit for fatal outcomes is 500–1,000 mg of ingested nicotine for an adult (6.5–13 mg/kg). Nicotine addiction involves drug-reinforced behavior, compulsive use, and relapse following abstinence. Nicotine dependence involves tolerance, sensitization, physical dependence, psychological dependence, and can cause distress. Nicotine withdrawal symptoms include depressed mood, stress, anxiety, irritability, difficulty concentrating, and sleep disturbances. Mild nicotine withdrawal symptoms are measurable in unrestricted smokers, who experience normal moods only as their blood nicotine levels peak, with each cigarette. On quitting, withdrawal symptoms worsen sharply, then gradually improve to a normal state. Nicotine use as a tool for quitting smoking has a good safety history. Animal studies suggest that nicotine may adversely affect cognitive development in adolescence, but the relevance of these findings to human brain development is disputed. At low amounts, it has a mild analgesic effect. According to the International Agency for Research on Cancer, "nicotine is not generally considered to be a carcinogen". The Surgeon General of the United States indicates that evidence is inadequate to infer the presence or absence of a causal relationship between exposure to nicotine and risk for cancer. Nicotine has been shown to produce birth defects in humans and is considered a teratogen. The median lethal dose of nicotine in humans is unknown. High doses are known to cause nicotine poisoning, organ failure, and death through paralysis of respiratory muscles, though serious or fatal overdoses are rare. Uses Medical The primary therapeutic use of nicotine is treating nicotine dependence to eliminate smoking and the damage it does to health. Controlled levels of nicotine are given to patients through gums, dermal patches, lozenges, inhalers, or nasal sprays to wean them off their dependence. A 2018 Cochrane Collaboration review found high-quality evidence that all current forms of nicotine replacement therapy (gum, patch, lozenges, inhaler, and nasal spray) increase the chances of successfully quitting smoking by , regardless of setting. Combining nicotine patch use with a faster acting nicotine replacement, like gum or spray, improves the odds of treatment success. In contrast to recreational nicotine products, which have been designed to maximize the likelihood of addiction, nicotine replacement products (NRTs) are designed to minimize addictiveness. The more quickly a dose of nicotine is delivered and absorbed, the higher the addiction risk. Investigative Nicotine is being researched in clinical trials for possible benefit in treating Parkinson's disease, dementia, attention deficit hyperactivity disorder (ADHD), and depression. Nicotine may partly attenuate sensory gating and attentional deficits associated with schizophrenia. Short-term use of transdermal nicotine was found to improve subjects’ reaction time and alertness in given tasks. Nicotine was not found to improve negative, positive, or other cognitive symptoms of schizophrenia. Pesticide Nicotine has been used as an insecticide since at least 1690, in the form of tobacco extracts or as pure nicotine sulphate (although other components of tobacco also seem to have pesticide effects). It acts on the nicotinic acetylcholine receptor, and gave the receptor its name. Nicotine is in IRAC group 4B. Nicotine insecticides have been banned in the US since 2014, including use on organic crops, and caution is recommended for small gardeners. Nicotine pesticides have been banned in the EU since 2009. Foods are imported from countries in which nicotine pesticides are allowed, such as China, but foods may not exceed maximum nicotine levels. Neonicotinoids, such as imidacloprid, which are derived from and structurally similar to nicotine, are widely used as agricultural and veterinary pesticides as of 2016. Performance Nicotine-containing products are sometimes used for the performance-enhancing effects of nicotine on cognition. A 2010 meta-analysis of 41 double-blind, placebo-controlled studies concluded that nicotine or smoking had significant positive effects on aspects of fine motor abilities, alerting and orienting attention, and episodic and working memory. A 2015 review noted that stimulation of the α4β2 nicotinic receptor is responsible for certain improvements in attentional performance; among the nicotinic receptor subtypes, nicotine has the highest binding affinity at the α4β2 receptor (ki=1 ), which is also the biological target that mediates nicotine's addictive properties. Nicotine has potential beneficial effects, but it also has paradoxical effects, which may be due to the inverted U-shape of the dose-response curve or pharmacokinetic features. Recreational Nicotine is used as a recreational drug. It is widely used, highly addictive and hard to discontinue. Nicotine is often used compulsively, and dependence can develop within days. Recreational drug users commonly use nicotine for its mood-altering effects. Recreational nicotine products include chewing tobacco, cigars, cigarettes, e-cigarettes, snuff, pipe tobacco, snus, and nicotine pouches. Alcohol infused with nicotine is called nicotini. Contraindications Nicotine use for tobacco cessation has few contraindications. It is not known whether nicotine replacement therapy is effective for smoking cessation in adolescents, as of 2014. It is therefore not recommended to adolescents. It is not safe to use nicotine during pregnancy or breastfeeding, although it is safer than smoking. The desirability of NRT use in pregnancy is therefore debated. Randomized trials and observational studies of nicotine replacement therapy in cardiovascular patients show no increase in adverse cardiovascular events compared to those treated with placebo. Using nicotine products during cancer treatment may be contraindicated, as nicotine may promote tumour growth, but temporary use of NRTs to quit smoking may be advised for harm reduction. Nicotine gum is contraindicated in individuals with temporomandibular joint disease. People with chronic nasal disorders and severe reactive airway disease require additional precautions when using nicotine nasal sprays. Nicotine in any form is contraindicated in individuals with a known hypersensitivity to nicotine. Adverse effects Nicotine is classified as a poison, and it is "extremely hazardous". However, at doses typically used by consumers, it presents little if any hazard to the user. A 2018 Cochrane Collaboration review lists nine main adverse events related to nicotine replacement therapy: headache, dizziness, lightheadedness, nausea, vomiting, gastrointestinal symptoms, insomnia, abnormal dreams, non-ischemic palpitations and chest pain, skin reactions, oral/nasal reactions, and hiccups. Many of these were also common in the placebo group without nicotine. Palpitations and chest pain were deemed "rare" and there was no evidence of an increased number of serious cardiac problems compared to the placebo group, even in people with established cardiac disease. The common side effects from nicotine exposure are listed in the table below. Serious adverse events due to the use of nicotine replacement therapy are extremely rare. At low amounts, it has a mild analgesic effect. However, at sufficiently high doses, nicotine may result in nausea, vomiting, diarrhea, salivation, bradycardia, and possibly seizures, hypoventilation, and death. Sleep Nicotine reduces the amount of rapid eye movement (REM) sleep, slow-wave sleep (SWS), and total sleep time in healthy nonsmokers given nicotine via a transdermal patch, and the reduction is dose-dependent. Acute nicotine intoxication has been found to significantly reduce total sleep time and increase REM latency, sleep onset latency, and non-rapid eye movement (NREM) stage 2 sleep time. Depressive non-smokers experience mood and sleep improvements under nicotine administration; however, subsequent nicotine withdrawal has a negative effect on both mood and sleep. Cardiovascular system Nicotine exerts several significant effects on the cardiovascular system. Primarily, it stimulates the sympathetic nervous system, leading to the release of catecholamines. This activation results in an increase in heart rate and blood pressure, as well as enhanced myocardial contractility, which raises the workload on the heart. Additionally, nicotine causes systemic vasoconstriction, including constriction of coronary arteries, which can reduce blood flow to the heart. Long-term exposure to nicotine may impair endothelial function, potentially contributing to atherosclerosis. Furthermore, nicotine has been associated with the development of cardiac arrhythmias, particularly in individuals who already have underlying heart disease. The effects of nicotine can be differentiated between short-term and long-term use. Short-term nicotine use, such as that associated with nicotine replacement therapy (NRT) for smoking cessation, appears to pose little cardiovascular risk, even for patients with known cardiovascular conditions. In contrast, longer-term nicotine use may not accelerate atherosclerosis but could contribute to acute cardiovascular events in those with pre-existing cardiovascular disease. Many severe cardiovascular effects traditionally associated with smoking may not be solely attributable to nicotine itself. Cigarette smoke contains numerous other potentially cardiotoxic substances, including carbon monoxide and oxidant gases. A 2016 review of the cardiovascular toxicity of nicotine concluded, "Based on current knowledge, we believe that the cardiovascular risks of nicotine from e-cigarette use in people without cardiovascular disease are quite low. We have concerns that nicotine from e-cigarettes could pose some risk for users with cardiovascular disease." A 2018 Cochrane review found that, in rare cases, nicotine replacement therapy can cause non-ischemic chest pain (i.e., chest pain that is unrelated to a heart attack) and heart palpitations, but does not increase the incidence of serious cardiac adverse events (i.e., myocardial infarction, stroke, and cardiac death) relative to controls. Blood pressure In the short term, nicotine causes a transient increase in blood pressure. Long term, epidemiological studies generally show increased blood pressure and hypertension among nicotine users. Reinforcement disorders Nicotine is highly addictive but paradoxically has quite weak reinforcing property compared to other drugs of abuse in various animals. Its addictiveness depends on how it is administered and also depends upon form in which nicotine is used. Animal research suggests that monoamine oxidase inhibitors, acetaldehyde and other constituents in tobacco smoke may enhance its addictiveness. Nicotine dependence involves aspects of both psychological dependence and physical dependence, since discontinuation of extended use has been shown to produce both affective (e.g., anxiety, irritability, craving, anhedonia) and somatic (mild motor dysfunctions such as tremor) withdrawal symptoms. Withdrawal symptoms peak in one to three days and can persist for several weeks. Even though other drugs of dependence can have withdrawal states lasting 6 months or longer, this does not appear to occur with cigarette withdrawal. Normal between-cigarettes discontinuation, in unrestricted smokers, causes mild but measurable nicotine withdrawal symptoms. These include mildly worse mood, stress, anxiety, cognition, and sleep, all of which briefly return to normal with the next cigarette. Smokers have a worse mood than they typically would have if they were not nicotine-dependent; they experience normal moods only immediately after smoking. Nicotine dependence is associated with poor sleep quality and shorter sleep duration among smokers. In dependent smokers, withdrawal causes impairments in memory and attention, and smoking during withdrawal returns these cognitive abilities to pre-withdrawal levels. The temporarily increased cognitive levels of smokers after inhaling smoke are offset by periods of cognitive decline during nicotine withdrawal. Therefore, the overall daily cognitive levels of smokers and non-smokers are roughly similar. Nicotine activates the mesolimbic pathway and induces long-term ΔFosB expression (i.e., produces phosphorylated ΔFosB isoforms) in the nucleus accumbens when inhaled or injected frequently or at high doses, but not necessarily when ingested. Consequently, high daily exposure (possibly excluding oral route) to nicotine can cause ΔFosB overexpression in the nucleus accumbens, resulting in nicotine addiction. Cancer Contrary to popular belief, nicotine itself does not cause cancer in humans, although it is unclear whether it functions as a tumor promoter . A 2018 report by the US National Academies of Sciences, Engineering, and Medicine concludes, "[w]hile it is biologically plausible that nicotine can act as a tumor promoter, the existing body of evidence indicates this is unlikely to translate into increased risk of human cancer." Although nicotine is classified as a non-carcinogenic substance, it can still promote tumor growth and metastasis. It induces several processes that contribute to cancer progression, including cell cycle progression, epithelial-to-mesenchymal transition, migration, invasion, angiogenesis, and evasion of apoptosis. These effects are primarily mediated through nicotinic acetylcholine receptors (nAChRs), particularly the α7 subtype, and to a lesser extent, β-adrenergic receptors (β-ARs). Activation of these receptors triggers several signaling cascades crucial in cancer biology, notably the MAPK/ERK pathway, PI3K/AKT pathway, and JAK-STAT signaling. Nicotine promotes lung cancer development by enhancing proliferation, angiogenesis, migration, invasion, and epithelial–mesenchymal transition (EMT) via nAChRs, which are present in lung cancer cells. Additionally, nicotine-induced EMT contributes to drug resistance in cancer cells. Nicotine in tobacco can form carcinogenic tobacco-specific nitrosamines through a nitrosation reaction. This occurs mostly in the curing and processing of tobacco. However, nicotine in the mouth and stomach can react to form N-nitrosonornicotine, a known type 1 carcinogen, suggesting that consumption of non-tobacco forms of nicotine may still play a role in carcinogenesis. Genotoxicity Nicotine causes DNA damage in several types of human cells as judged by assays for genotoxicity such as the comet assay, cytokinesis-block micronucleus test and chromosome aberrations test. In humans, this damage can happen in primary parotid gland cells, lymphocytes, and respiratory tract cells. Pregnancy and breastfeeding Nicotine has been shown to produce birth defects in some animal species, but not others; consequently, it is considered to be a possible teratogen in humans. In animal studies that resulted in birth defects, researchers found that nicotine negatively affects fetal brain development and pregnancy outcomes; the negative effects on early brain development are associated with abnormalities in brain metabolism and neurotransmitter system function. Nicotine crosses the placenta and is found in the breast milk of mothers who smoke as well as mothers who inhale passive smoke. Nicotine exposure in utero is responsible for several complications of pregnancy and birth: pregnant women who smoke are at greater risk for both miscarriage and stillbirth and infants exposed to nicotine in utero tend to have lower birth weights. A McMaster University research group observed in 2010 that rats exposed to nicotine in the womb (via parenteral infusion) later in life had conditions including type 2 diabetes, obesity, hypertension, neurobehavioral defects, respiratory dysfunction, and infertility. Overdose It is unlikely that a person would overdose on nicotine through smoking alone. The US Food and Drug Administration (FDA) stated in 2013 that there are no significant safety concerns associated with the use of more than one form of over-the-counter (OTC) nicotine replacement therapy at the same time, or using OTC NRT at the same time as another nicotine-containing product, like cigarettes. The median lethal dose of nicotine in humans is unknown. Nevertheless, nicotine has a relatively high toxicity in comparison to many other alkaloids such as caffeine, which has an LD50 of 127 mg/kg when administered to mice. At sufficiently high doses, it is associated with nicotine poisoning, which, while common in children (in whom poisonous and lethal levels occur at lower doses per kilogram of body weight) rarely results in significant morbidity or death. The estimated lower dose limit for fatal outcomes is 500–1,000 mg of ingested nicotine for an adult (6.5–13 mg/kg). The initial symptoms of a nicotine overdose typically include nausea, vomiting, diarrhea, hypersalivation, abdominal pain, tachycardia (rapid heart rate), hypertension (high blood pressure), tachypnea (rapid breathing), headache, dizziness, pallor (pale skin), auditory or visual disturbances, and perspiration, followed shortly after by marked bradycardia (slow heart rate), bradypnea (slow breathing), and hypotension (low blood pressure). An increased respiratory rate (i.e., tachypnea) is one of the primary signs of nicotine poisoning. At sufficiently high doses, somnolence (sleepiness or drowsiness), confusion, syncope (loss of consciousness from fainting), shortness of breath, marked weakness, seizures, and coma may occur. Lethal nicotine poisoning rapidly produces seizures, and death – which may occur within minutes – is believed to be due to respiratory paralysis. Toxicity Today nicotine is less commonly used in agricultural insecticides, which was a main source of poisoning. More recent cases of poisoning typically appear to be in the form of Green Tobacco Sickness (GTS), accidental ingestion of tobacco or tobacco products, or ingestion of nicotine-containing plants. People who harvest or cultivate tobacco may experience GTS, a type of nicotine poisoning caused by dermal exposure to wet tobacco leaves. This occurs most commonly in young, inexperienced tobacco harvesters who do not consume tobacco. People can be exposed to nicotine in the workplace by breathing it in, skin absorption, swallowing it, or eye contact. The Occupational Safety and Health Administration (OSHA) has set the legal limit (permissible exposure limit) for nicotine exposure in the workplace as 0.5 mg/m3 skin exposure over an 8-hour workday. The US National Institute for Occupational Safety and Health (NIOSH) has set a recommended exposure limit (REL) of 0.5 mg/m3 skin exposure over an 8-hour workday. At environmental levels of 5 mg/m3, nicotine is immediately dangerous to life and health. Drug interactions Pharmacodynamic Potential interaction with sympathomimetic drugs (adrenergic agonists) and sympatholytic drugs (alpha-blockers and beta-blockers). Pharmacokinetic Nicotine and cigarette smoke both induce the expression of liver enzymes (e.g., certain cytochrome P450 proteins) which metabolize drugs, leading to the potential for alterations in drug metabolism. Smoking cessation may decrease the metabolism of acetaminophen, beta-blockers, caffeine, oxazepam, pentazocine, propoxyphene, theophylline, and tricyclic antidepressants, leading to higher plasma concentrations of these drugs. Possible alteration of nicotine absorption through the skin from the transdermal nicotine patch by drugs that cause vasodilation or vasoconstriction. Possible alteration of nicotine absorption through the nasal cavity from the nicotine nasal spray by nasal vasoconstrictors (e.g., xylometazoline). Possible alteration of nicotine absorption through oral mucosa from nicotine gum and lozenges by food and drink that modify salivary pH. Pharmacology Pharmacodynamics Nicotine acts as a receptor agonist at most nicotinic acetylcholine receptors (nAChRs), except at two nicotinic receptor subunits (nAChRα9 and nAChRα10) where it acts as a receptor antagonist. Such antagonism results in mild analgesia. Central nervous system By binding to nicotinic acetylcholine receptors in the brain, nicotine elicits its psychoactive effects and increases the levels of several neurotransmitters in various brain structures – acting as a sort of "volume control". Nicotine has a higher affinity for nicotinic receptors in the brain than those in skeletal muscle, though at toxic doses it can induce contractions and respiratory paralysis. Nicotine's selectivity is thought to be due to a particular amino acid difference on these receptor subtypes. Nicotine is unusual in comparison to most drugs, as its profile changes from stimulant to sedative with increasing dosages, a phenomenon known as "Nesbitt's paradox" after the doctor who first described it in 1969. At very high doses it dampens neuronal activity. Nicotine induces both behavioral stimulation and anxiety in animals. Research into nicotine's most predominant metabolite, cotinine, suggests that some of nicotine's psychoactive effects are mediated by cotinine. Nicotine activates nicotinic receptors (particularly α4β2 nicotinic receptors, but also α5 nAChRs) on neurons that innervate the ventral tegmental area and within the mesolimbic pathway where it appears to cause the release of dopamine. This nicotine-induced dopamine release occurs at least partially through activation of the cholinergic–dopaminergic reward link in the ventral tegmental area. Nicotine can modulate the firing rate of the ventral tegmental area neurons. These actions are largely responsible for the strongly reinforcing effects of nicotine, which often occur in the absence of euphoria; however, mild euphoria from nicotine use can occur in some individuals. Chronic nicotine use inhibits class I and II histone deacetylases in the striatum, where this effect plays a role in nicotine addiction. Sympathetic nervous system Nicotine also activates the sympathetic nervous system, acting via splanchnic nerves to the adrenal medulla, stimulating the release of epinephrine. Acetylcholine released by preganglionic sympathetic fibers of these nerves acts on nicotinic acetylcholine receptors, causing the release of epinephrine (and norepinephrine) into the bloodstream. Adrenal medulla By binding to ganglion type nicotinic receptors in the adrenal medulla, nicotine increases flow of adrenaline (epinephrine), a stimulating hormone and neurotransmitter. By binding to the receptors, it causes cell depolarization and an influx of calcium through voltage-gated calcium channels. Calcium triggers the exocytosis of chromaffin granules and thus the release of epinephrine (and norepinephrine) into the bloodstream. The release of epinephrine (adrenaline) causes an increase in heart rate, blood pressure and respiration, as well as higher blood glucose levels. Pharmacokinetics As nicotine enters the body, it is distributed quickly through the bloodstream and crosses the blood–brain barrier reaching the brain within 10–20 seconds after inhalation. The elimination half-life of nicotine in the body is around two hours. Nicotine is primarily excreted in urine and urinary concentrations vary depending upon urine flow rate and urine pH. The amount of nicotine absorbed by the body from smoking can depend on many factors, including the types of tobacco, whether the smoke is inhaled, and whether a filter is used. However, it has been found that the nicotine yield of individual products has only a small effect (4.4%) on the blood concentration of nicotine, suggesting "the assumed health advantage of switching to lower-tar and lower-nicotine cigarettes may be largely offset by the tendency of smokers to compensate by increasing inhalation". Nicotine has a half-life of 1–2 hours. Cotinine is an active metabolite of nicotine that remains in the blood with a half-life of 18–20 hours, making it easier to analyze. Nicotine is metabolized in the liver by cytochrome P450 enzymes (mostly CYP2A6, and also by CYP2B6) and FMO3, which selectively metabolizes (S)-nicotine. A major metabolite is cotinine. Other primary metabolites include nicotine N-oxide, nornicotine, nicotine isomethonium ion, 2-hydroxynicotine and nicotine glucuronide. Under some conditions, other substances may be formed such as myosmine. Glucuronidation and oxidative metabolism of nicotine to cotinine are both inhibited by menthol, an additive to mentholated cigarettes, thus increasing the half-life of nicotine in vivo. Metabolism Nicotine decreases hunger and as a consequence food consumption, alongside increasing energy expenditure. The majority of research shows that nicotine reduces body weight, but some researchers have found that nicotine may result in weight gain under specific types of eating habits in animal models. Nicotine effect on weight appears to result from nicotine's stimulation of α3β4 nAChR receptors located in the POMC neurons in the arcuate nucleus and subsequently the melanocortin system, especially the melanocortin-4 receptors on second-order neurons in the paraventricular nucleus of the hypothalamus, thus modulating feeding inhibition. POMC neurons are a precursor of the melanocortin system, a critical regulator of body weight and peripheral tissue such as skin and hair. Chemistry Nicotine is a hygroscopic, colorless to yellow-brown, oily liquid, that is readily soluble in alcohol, ether or light petroleum. It is miscible with water in its neutral amine base form between 60 °C and 210 °C. It is a dibasic nitrogenous base, having Kb1=1×10−6, Kb2=1×10−11. It readily forms ammonium salts with acids that are usually solid and water-soluble. Its flash point is 95 °C and its auto-ignition temperature is 244 °C. Nicotine is readily volatile (vapor pressure 5.5 Pa at 25 °C) On exposure to ultraviolet light or various oxidizing agents, nicotine is converted to nicotine oxide, nicotinic acid (niacin, vitamin B3), and methylamine. Nicotine is chiral and hence optically active, having two enantiomeric forms. The naturally occurring form of nicotine is levorotatory with a specific rotation of [α]D=–166.4° ((−)-nicotine). The dextrorotatory form, (+)-nicotine is physiologically less active than (−)-nicotine. (−)-nicotine is more toxic than (+)-nicotine. The salts of (−)-nicotine are usually dextrorotatory; this conversion between levorotatory and dextrorotatory upon protonation is common among alkaloids. The hydrochloride and sulfate salts become optically inactive if heated in a closed vessel above 180 °C. Anabasine is a structural isomer of nicotine, as both compounds have the molecular formula . Nicotine that is found in natural tobacco is primarily (99%) the S-enantiomer. Conversely, the most common chemistry synthetic methods for generating nicotine yields a product that is approximately equal proportions of the S- and R-enantiomers. This suggests that tobacco-derived and synthetic nicotine can be determined by measuring the ratio of the two different enantiomers, although means exist for adjusting the relative levels of the enantiomers or performing a synthesis that only leads to the S-enantiomer. There is limited data on the relative physiological effects of these two enantiomers, especially in people. However, the studies to date indicate that (S)-nicotine is more potent than (R)-nicotine and (S)-nicotine causes stronger sensations or irritation than (R)-nicotine. Studies have not been adequate to determine the relative addictiveness of the two enantiomers in people. Pod mod electronic cigarettes use nicotine in the form of a protonated nicotine, rather than free-base nicotine found in earlier generations. Preparation The first laboratory preparation of nicotine (as its racemate) was described in 1904. The starting material was an N-substituted pyrrole derivative, which was heated to convert it by a [1,5] sigmatropic shift to the isomer with a carbon bond between the pyrrole and pyridine rings, followed by methylation and selective reduction of the pyrrole ring using tin and hydrochloric acid. Many other syntheses of nicotine, in both racemic and chiral forms have since been published. Biosynthesis The biosynthetic pathway of nicotine involves a coupling reaction between the two cyclic structures that comprise nicotine. Metabolic studies show that the pyridine ring of nicotine is derived from niacin (nicotinic acid) while the pyrrolidine is derived from N-methyl-Δ1-pyrrollidium cation. Biosynthesis of the two component structures proceeds via two independent syntheses, the NAD pathway for niacin and the tropane pathway for N-methyl-Δ1-pyrrollidium cation. The NAD pathway in the genus Nicotiana begins with the oxidation of aspartic acid into α-amino succinate by aspartate oxidase (AO). This is followed by a condensation with glyceraldehyde-3-phosphate and a cyclization catalyzed by quinolinate synthase (QS) to give quinolinic acid. Quinolinic acid then reacts with phosphoribosyl pyrophosphate catalyzed by quinolinic acid phosphoribosyl transferase (QPT) to form niacin mononucleotide (NaMN). The reaction now proceeds via the NAD salvage cycle to produce niacin via the conversion of nicotinamide by the enzyme nicotinamidase. The N-methyl-Δ1-pyrrollidium cation used in the synthesis of nicotine is an intermediate in the synthesis of tropane-derived alkaloids. Biosynthesis begins with decarboxylation of ornithine by ornithine decarboxylase (ODC) to produce putrescine. Putrescine is then converted into N-methyl putrescine via methylation by SAM catalyzed by putrescine N-methyltransferase (PMT). N-methyl putrescine then undergoes deamination into 4-methylaminobutanal by the N-methyl putrescine oxidase (MPO) enzyme, 4-methylaminobutanal then spontaneously cyclize into N-methyl-Δ1-pyrrollidium cation. The final step in the synthesis of nicotine is the coupling between N-methyl-Δ1-pyrrollidium cation and niacin. Although studies conclude some form of coupling between the two component structures, the definite process and mechanism remains undetermined. The current agreed theory involves the conversion of niacin into 2,5-dihydropyridine through 3,6-dihydronicotinic acid. The 2,5-dihydropyridine intermediate would then react with N-methyl-Δ1-pyrrollidium cation to form enantiomerically pure (−)-nicotine. Detection in body fluids Nicotine can be quantified in blood, plasma, or urine to confirm a diagnosis of poisoning or to facilitate a medicolegal death investigation. Urinary or salivary cotinine concentrations are frequently measured for the purposes of pre-employment and health insurance medical screening programs. Careful interpretation of results is important, since passive exposure to cigarette smoke can result in significant accumulation of nicotine, followed by the appearance of its metabolites in various body fluids. Nicotine use is not regulated in competitive sports programs. Methods for analysis of enantiomers Methods for measuring the two enantiomers are straightforward and include normal-phase liquid chromatography, liquid chromatography with a chiral column. However, since methods can be used to alter the two enantiomers, it may not be possible to distinguish tobacco-derived from synthetic nicotine simply by measuring the levels of the two enantiomers. A new approach uses hydrogen and deuterium nuclear magnetic resonance to distinguish tobacco-derived and synthetic nicotine based on differences the substrates used in the natural synthetic pathway performed in the tobacco plant and the substrates most used in synthesis. Another approach measures the carbon-14 content which also differs between natural and laboratory-based tobacco. These methods remain to be fully evaluated and validated using a wide range of samples. Natural occurrence Nicotine is a secondary metabolite produced in a variety of plants in the family Solanaceae, most notably in tobacco Nicotiana tabacum, where it can be found at high concentrations of 0.5 to 7.5%. Nicotine is also found in the leaves of other tobacco species, such as Nicotiana rustica (in amounts of 2–14%). Nicotine production is strongly induced in response to wounding as part of a jasmonate-dependent reaction. Specialist insects on tobacco, such as the tobacco hornworm (Manduca sexta), have a number of adaptations to the detoxification and even adaptive re-purposing of nicotine. Nicotine is also found at low concentrations in the nectar of tobacco plants, where it may promote outcrossing by affecting the behavior of hummingbird pollinators. Nicotine occurs in smaller amounts (varying from 2–7 μg/kg, or 20–70 millionths of a percent wet weight) in other Solanaceaeous plants, including some crop species such as potatoes, tomatoes, eggplant, and peppers, as well as non-crop species such as Duboisia hopwoodii. The amounts of nicotine in tomatoes lowers substantially as the fruit ripens. A 1999 report found "In some papers it is suggested that the contribution of dietary nicotine intake is significant when compared with exposure to ETS [environmental tobacco smoke] or by active smoking of small numbers of cigarettes. Others consider the dietary intake to be negligible unless inordinately large amounts of specific vegetables are consumed." The amount of nicotine eaten per day is roughly around 1.4 and 2.25 μg/day at the 95th percentile. These numbers may be low due to insufficient food intake data. The concentrations of nicotine in vegetables are difficult to measure accurately, since they are very low (parts per billion range). Pure nicotine tastes "terrible". History, society and culture Nicotine was originally isolated from the tobacco plant in 1828 by chemists Wilhelm Heinrich Posselt and Karl Ludwig Reimann from Germany, who believed it was a poison. Its chemical empirical formula was described by Melsens in 1843, its structure was discovered by Adolf Pinner and Richard Wolffenstein in 1893, and it was first synthesized by Amé Pictet and A. Rotschy in 1904. Nicotine is named after the tobacco plant Nicotiana tabacum, which in turn is named after the French ambassador in Portugal, Jean Nicot de Villemain, who sent tobacco and seeds to Paris in 1560, presented to the French King, and who promoted their medicinal use. Smoking was believed to protect against illness, particularly the plague. Tobacco was introduced to Europe in 1559, and by the late 17th century, it was used not only for smoking but also as an insecticide. After World War II, over 2,500 tons of nicotine insecticide were used worldwide, but by the 1980s the use of nicotine insecticide had declined below 200 tons. This was due to the availability of other insecticides that are cheaper and less harmful to mammals. The nicotine content of popular American-brand cigarettes has increased over time, and one study found that there was an average increase of 1.78% per year between the years of 1998 and 2005. Although methods of production of synthetic nicotine have existed for decades, it was believed that the cost of making nicotine by laboratory synthesis was cost prohibitive compared to extracting nicotine from tobacco. However, recently synthetic nicotine started to be found in different brands of e-cigarettes and oral pouches and marketed as "tobacco-free." The US FDA is tasked with reviewing tobacco products such as e-cigarettes and determining which can be authorized for sale. In response to the likelihood that FDA would not authorize many e-cigarettes to be marketed, e-cigarette companies began marketing products that they claimed to contain nicotine that were not made or derived from tobacco, but contained synthetic nicotine instead, and thus, would be outside FDA's tobacco regulatory authority. Similarly, nicotine pouches that claimed to contain non-tobacco (synthetic) nicotine were also introduced. The cost of synthetic nicotine has decreased as the market for the product increased. In March 2022, the U.S. Congress passed a law (the Consolidated Appropriations Act, 2022) that expanded FDA's tobacco regulatory authority to include tobacco products containing nicotine from any source, thereby including products made with synthetic nicotine. Legal status In the United States, nicotine products and nicotine replacement therapy products like Nicotrol are only available to people 18 and above; proof of age is required; not for sale in vending machine or from any source where proof of age cannot be verified. As of 2019, the minimum age to purchase tobacco in the US is 21 at the federal level. In the European Union, the minimum age to purchase nicotine products is 18. However, there is no minimum age requirement to use tobacco or nicotine products. In the United Kingdom, the Tobacco and Related Products Regulations 2016 implemented the European directive 2014/40/EU, amended by Tobacco Products and Nicotine Inhaling Products (Amendment etc.) (EU Exit) Regulations 2019 and the Tobacco Products and Nicotine Inhaling Products (Amendment) (EU Exit) Regulations 2020. Additionally other regulations limit advertising, sale and display of tobacco products and other products containing nicotine for human consumption. The Sunak government proposed banning disposable vapes to limit their appeal and affordability for children and to reduce the amount of waste generated. In media In some anti-smoking literature, the harm that tobacco smoking and nicotine addiction does is personified as Nick O'Teen, represented as a humanoid with some aspect of a cigarette or cigarette butt about him or his clothes and hat. Nick O'Teen was a villain that was created for the Health Education Council. The character was featured in three animated anti-smoking public service announcements in which he tries to get kids addicted to cigarettes before being foiled by the DC Comics character Superman. Nicotine was often compared to caffeine in advertisements in the 1980s by the tobacco industry, and later in the 2010s by the electronic cigarettes industry, in an effort to reduce the stigmatization and the public perception of the risks associated with nicotine use. Research Central nervous system While acute/initial nicotine intake causes activation of neuronal nicotine receptors, chronic low doses of nicotine use leads to desensitization of those receptors (due to the development of tolerance) and results in an antidepressant effect, with early research showing low dose nicotine patches could be an effective treatment of major depressive disorder in non-smokers. Though tobacco smoking is associated with an increased risk of Alzheimer's disease, there is evidence that nicotine itself has the potential to prevent and treat Alzheimer's disease. Smoking is associated with a decreased risk of Parkinson's disease; however, it is unknown whether this is due to people with healthier brain dopaminergic reward centers (the area of the brain affected by Parkinson's) being more likely to enjoy smoking and thus pick up the habit, nicotine directly acting as a neuroprotective agent, or other compounds in cigarette smoke acting as neuroprotective agents. Immune system Immune cells of both the innate immune system and adaptive immune systems frequently express the α2, α5, α6, α7, α9, and α10 subunits of nicotinic acetylcholine receptors. Evidence suggests that nicotinic receptors which contain these subunits are involved in the regulation of immune function. Optopharmacology A photoactivatable form of nicotine, which releases nicotine when exposed to ultraviolet light with certain conditions, has been developed for studying nicotinic acetylcholine receptors in brain tissue. Oral health Several in vitro studies have investigated the potential effects of nicotine on a range of oral cells. A recent systematic review concluded that nicotine was unlikely to be cytotoxic to oral cells in vitro in most physiological conditions but further research is needed. Understanding the potential role of nicotine in oral health has become increasingly important given the recent introduction of novel nicotine products and their potential role in helping smokers quit.	Biology and health sciences	Biochemistry and molecular biology	null
38300	https://en.wikipedia.org/wiki/Pancreas	Pancreas	The pancreas is an organ of the digestive system and endocrine system of vertebrates. In humans, it is located in the abdomen behind the stomach and functions as a gland. The pancreas is a mixed or heterocrine gland, i.e., it has both an endocrine and a digestive exocrine function. 99% of the pancreas is exocrine and 1% is endocrine. As an endocrine gland, it functions mostly to regulate blood sugar levels, secreting the hormones insulin, glucagon, somatostatin and pancreatic polypeptide. As a part of the digestive system, it functions as an exocrine gland secreting pancreatic juice into the duodenum through the pancreatic duct. This juice contains bicarbonate, which neutralizes acid entering the duodenum from the stomach; and digestive enzymes, which break down carbohydrates, proteins and fats in food entering the duodenum from the stomach. Inflammation of the pancreas is known as pancreatitis, with common causes including chronic alcohol use and gallstones. Because of its role in the regulation of blood sugar, the pancreas is also a key organ in diabetes mellitus. Pancreatic cancer can arise following chronic pancreatitis or due to other reasons, and carries a very poor prognosis, as it is often only identified after it has spread to other areas of the body. The word pancreas comes from the Greek πᾶν (pân, "all") & κρέας (kréas, "flesh"). The function of the pancreas in diabetes has been known since at least 1889, with its role in insulin production identified in 1921. Structure The pancreas is an organ that in humans lies in the abdomen, stretching from behind the stomach to the left upper abdomen near the spleen. In adults, it is about long, , and salmon-coloured in appearance. Anatomically, the pancreas is divided into a head, neck, body, and tail. The pancreas stretches from the inner curvature of the duodenum, where the head surrounds two blood vessels: the superior mesenteric artery and vein. The longest part of the pancreas, the body, stretches across behind the stomach, and the tail of the pancreas ends adjacent to the spleen. Two ducts, the main pancreatic duct and a smaller accessory pancreatic duct run through the body of the pancreas. The main pancreatic duct joins with the common bile duct forming a small ballooning called the ampulla of Vater (hepatopancreatic ampulla). This ampulla is surrounded by a muscle, the sphincter of Oddi. This ampulla opens into the descending part of the duodenum. The opening of the common bile duct into main pancreatic duct is controlled by sphincter of Boyden. The accessory pancreatic duct opens into duodenum with separate openings located above the opening of the main pancreatic duct. Parts The head of the pancreas sits within the curvature of the duodenum, and wraps around the superior mesenteric artery and vein. To the right sits the descending part of the duodenum, and between these travel the superior and inferior pancreaticoduodenal arteries. Behind rest the inferior vena cava, and the common bile duct. In front sit the peritoneal membrane and the transverse colon. A small uncinate process emerges from below the head, situated behind the superior mesenteric vein and sometimes artery. The neck of the pancreas separates the head of the pancreas, located in the curvature of the duodenum, from the body. The neck is about wide, and sits in front of where the portal vein is formed. The neck lies mostly behind the pylorus of the stomach, and is covered with peritoneum. The anterior superior pancreaticoduodenal artery travels in front of the neck of the pancreas. The body is the largest part of the pancreas, and mostly lies behind the stomach, tapering along its length. The peritoneum sits on top of the body of the pancreas, and the transverse colon in front of the peritoneum. Behind the pancreas are several blood vessels, including the aorta, the splenic vein, and the left renal vein, as well as the beginning of the superior mesenteric artery. Below the body of the pancreas sits some of the small intestine, specifically the last part of the duodenum and the jejunum to which it connects, as well as the suspensory ligament of the duodenum which falls between these two. In front of the pancreas sits the transverse colon. The pancreas narrows towards the tail, which sits near to the spleen. It is usually between long, and sits between the layers of the ligament between the spleen and the left kidney. The splenic artery and vein, which also passes behind the body of the pancreas, pass behind the tail of the pancreas. Blood supply The pancreas has a rich blood supply, with vessels originating as branches of both the coeliac artery and superior mesenteric artery. The splenic artery, the largest branch of the celiac trunk, runs along the top of the pancreas, and supplies the left part of the body and the tail of the pancreas through its pancreatic branches, the largest of which is called the greater pancreatic artery. The superior and inferior pancreaticoduodenal arteries run along the back and front surfaces of the head of the pancreas adjacent to the duodenum. These supply the head of the pancreas. These vessels join together () in the middle. The body and neck of the pancreas drain into the splenic vein, which sits behind the pancreas. The head drains into, and wraps around, the superior mesenteric and portal veins, via the pancreaticoduodenal veins. The pancreas drains into lymphatic vessels that travel alongside its arteries, and has a rich lymphatic supply. The lymphatic vessels of the body and tail drain into splenic lymph nodes, and eventually into lymph nodes that lie in front of the aorta, between the coeliac and superior mesenteric arteries. The lymphatic vessels of the head and neck drain into intermediate lymphatic vessels around the pancreaticoduodenal, mesenteric and hepatic arteries, and from there into the lymph nodes that lie in front of the aorta. Microanatomy The pancreas contains tissue with an endocrine and exocrine role, and this division is also visible when the pancreas is viewed under a microscope. The majority of pancreatic tissue has a digestive role. The cells with this role form clusters () around small ducts, and are arranged in lobes that have thin fibrous walls. The cells of each acinus secrete inactive digestive enzymes called zymogens into the small intercalated ducts which they surround. In each acinus, the cells are pyramid-shaped and situated around the intercalated ducts, with the nuclei resting on the basement membrane, a large endoplasmic reticulum, and a number of zymogen granules visible within the cytoplasm. The intercalated ducts drain into larger intralobular ducts within the lobule, and finally interlobular ducts. The ducts are lined by a single layer of column-shaped cells. There is more than one layer of cells as the diameter of the ducts increases. The tissues with an endocrine role within the pancreas exist as clusters of cells called pancreatic islets (also called islets of Langerhans) that are distributed throughout the pancreas. Pancreatic islets contain alpha cells, beta cells, and delta cells, each of which releases a different hormone. These cells have characteristic positions, with alpha cells (secreting glucagon) tending to be situated around the periphery of the islet, and beta cells (secreting insulin) more numerous and found throughout the islet. Enterochromaffin cells are also scattered throughout the islets. Islets are composed of up to 3,000 secretory cells, and contain several small arterioles to receive blood, and venules that allow the hormones secreted by the cells to enter the systemic circulation. Variation The size of the pancreas varies considerably. Several anatomical variations exist, relating to the embryological development of the two pancreatic buds. The pancreas develops from these buds on either side of the duodenum. The ventral bud rotates to lie next to the dorsal bud, eventually fusing. In about 10% of adults, an accessory pancreatic duct may be present if the main duct of the dorsal bud of the pancreas does not regress; this duct opens into the minor duodenal papilla. If the two buds themselves, each having a duct, do not fuse, a pancreas may exist with two separate ducts, a condition known as a pancreas divisum. This condition has no physiologic consequence. If the ventral bud does not fully rotate, an annular pancreas may exist, where part or all of the duodenum is encircled by the pancreas. This may be associated with duodenal atresia. Gene and protein expression 10,000 protein coding genes (~50% of all human genes) are expressed in the normal human pancreas. Less than 100 of these genes are specifically expressed in the pancreas. Similar to the salivary glands, most pancreas-specific genes encode for secreted proteins. Corresponding pancreas-specific proteins are either expressed in the exocrine cellular compartment and have functions related to digestion or food uptake such as digestive chymotrypsinogen enzymes and pancreatic lipase PNLIP, or are expressed in the various cells of the endocrine pancreatic islets and have functions related to secreted hormones such as insulin, glucagon, somatostatin and pancreatic polypeptide. Development The pancreas forms during development from two buds that arise from the duodenal part of the foregut, an embryonic tube that is a precursor to the gastrointestinal tract. It is of endodermal origin. Pancreatic development begins with the formation of a dorsal and ventral pancreatic bud. Each joins with the foregut through a duct. The dorsal pancreatic bud forms the neck, body, and tail of the developed pancreas, and the ventral pancreatic bud forms the head and uncinate process. The definitive pancreas results from rotation of the ventral bud and the fusion of the two buds. During development, the duodenum rotates to the right, and the ventral bud rotates with it, moving to a position that becomes more dorsal. Upon reaching its final destination, the ventral pancreatic bud is below the larger dorsal bud, and eventually fuses with it. At this point of fusion, the main ducts of the ventral and dorsal pancreatic buds fuse, forming the main pancreatic duct. Usually, the duct of the dorsal bud regresses, leaving the main pancreatic duct. Cellular development Pancreatic progenitor cells are precursor cells that differentiate into the functional pancreatic cells, including exocrine acinar cells, endocrine islet cells, and ductal cells. These progenitor cells are characterised by the co-expression of the transcription factors PDX1 and NKX6-1. The cells of the exocrine pancreas differentiate through molecules that induce differentiation including follistatin, fibroblast growth factors, and activation of the Notch receptor system. Development of the exocrine acini progresses through three successive stages. These are the predifferentiated, protodifferentiated, and differentiated stages, which correspond to undetectable, low, and high levels of digestive enzyme activity, respectively. Pancreatic progenitor cells differentiate into endocrine islet cells under the influence of neurogenin-3 and ISL1, but only in the absence of notch receptor signaling. Under the direction of a Pax gene, the endocrine precursor cells differentiate to form alpha and gamma cells. Under the direction of Pax-6, the endocrine precursor cells differentiate to form beta and delta cells. The pancreatic islets form as the endocrine cells migrate from the duct system to form small clusters around capillaries. This occurs around the third month of development, and insulin and glucagon can be detected in the human fetal circulation by the fourth or fifth month of development. Function The pancreas is involved in blood sugar control and metabolism within the body, and also in the secretion of substances (collectively pancreatic juice) that help digestion. These are divided into an "endocrine" role, relating to the secretion of insulin and other substances within pancreatic islets that help control blood sugar levels and metabolism within the body, and an "exocrine" role, relating to the secretion of enzymes involved in digesting substances in the digestive tract. Blood glucose regulation Cells within the pancreas help to maintain blood glucose levels (homeostasis). The cells that do this are located within the pancreatic islets that are present throughout the pancreas. When blood glucose levels are low, alpha cells secrete glucagon, which increases blood glucose levels. When blood glucose levels are high beta cells secrete insulin to decrease glucose in blood. Delta cells in the islet also secrete somatostatin which decreases the release of insulin and glucagon. Glucagon acts to increase glucose levels by promoting the creation of glucose and the breakdown of glycogen to glucose in the liver. It also decreases the uptake of glucose in fat and muscle. Glucagon release is stimulated by low blood glucose or insulin levels, and during exercise. Insulin acts to decrease blood glucose levels by facilitating uptake by cells (particularly skeletal muscle), and promoting its use in the creation of proteins, fats and carbohydrates. Insulin is initially created as a precursor form called preproinsulin. This is converted to proinsulin and cleaved by C-peptide to insulin which is then stored in granules in beta cells. Glucose is taken into the beta cells and degraded. The end effect of this is to cause depolarisation of the cell membrane which stimulates the release of the insulin. The main factor influencing the secretion of insulin and glucagon are the levels of glucose in blood plasma. Low blood sugar stimulates glucagon release, and high blood sugar stimulates insulin release. Other factors also influence the secretion of these hormones. Some amino acids, that are byproducts of the digestion of protein, stimulate insulin and glucagon release. Somatostatin acts as an inhibitor of both insulin and glucagon. The autonomic nervous system also plays a role. Activation of Beta-2 receptors of the sympathetic nervous system by catecholamines secreted from sympathetic nerves stimulates secretion of insulin and glucagon, whereas activation of Alpha-1 receptors inhibits secretion. M3 receptors of the parasympathetic nervous system act when stimulated by the right vagus nerve to stimulate release of insulin from beta cells. Digestion The pancreas plays a vital role in the digestive system. It does this by secreting a fluid that contains digestive enzymes into the duodenum, the first part of the small intestine that receives food from the stomach. These enzymes help to break down carbohydrates, proteins and lipids (fats). This role is called the "exocrine" role of the pancreas. The cells that do this are arranged in clusters called acini. Secretions into the middle of the acinus accumulate in intralobular ducts, which drain to the main pancreatic duct, which drains directly into the duodenum. About 1.5–3 liters of fluid are secreted in this manner every day. The cells in each acinus are filled with granules containing the digestive enzymes. These are secreted in an inactive form termed zymogens or proenzymes. When released into the duodenum, they are activated by the enzyme enterokinase present in the lining of the duodenum. The proenzymes are cleaved, creating a cascade of activating enzymes. Enzymes that break down proteins begin with activation of trypsinogen to trypsin. The free trypsin then cleaves the rest of the trypsinogen, as well as chymotrypsinogen to its active form chymotrypsin. Enzymes secreted involved in the digestion of fats include lipase, phospholipase A2, lysophospholipase, and cholesterol esterase. Enzymes that break down starch and other carbohydrates include amylase. These enzymes are secreted in a fluid rich in bicarbonate. Bicarbonate helps maintain an alkaline pH for the fluid, a pH in which most of the enzymes act most efficiently, and also helps to neutralise the stomach acids that enter the duodenum. Secretion is influenced by hormones including secretin, cholecystokinin, and VIP, as well as acetylcholine stimulation from the vagus nerve. Secretin is released from the S cells which form part of the lining of the duodenum in response to stimulation by gastric acid. Along with VIP, it increases the secretion of enzymes and bicarbonate. Cholecystokinin is released from Ito cells of the lining of the duodenum and jejunum mostly in response to long chain fatty acids, and increases the effects of secretin. At a cellular level, bicarbonate is secreted from centroacinar and ductal cells through a sodium and bicarbonate cotransporter that acts because of membrane depolarisation caused by the cystic fibrosis transmembrane conductance regulator. Secretin and VIP act to increase the opening of the cystic fibrosis transmembrane conductance regulator, which leads to more membrane depolarisation and more secretion of bicarbonate. A variety of mechanisms act to ensure that the digestive action of the pancreas does not act to digest pancreatic tissue itself. These include the secretion of inactive enzymes (zymogens), the secretion of the protective enzyme trypsin inhibitor, which inactivates trypsin, the changes in pH that occur with bicarbonate secretion that stimulate digestion only when the pancreas is stimulated, and the fact that the low calcium within cells causes inactivation of trypsin. Additional functions The pancreas also secretes vasoactive intestinal peptide and pancreatic polypeptide. Enterochromaffin cells of the pancreas secrete the hormones motilin, serotonin, and substance P. It has been demonstrated that pancreatic tissue is a strong accumulator and secretor in the intestine of radioactive cesium (Cs-137). Clinical significance Inflammation Inflammation of the pancreas is known as pancreatitis. Pancreatitis is most often associated with recurrent gallstones or chronic alcohol use, with other common causes including traumatic damage, damage following an ERCP, some medications, infections such as mumps and very high blood triglyceride levels. Acute pancreatitis is likely to cause intense pain in the central abdomen, that often radiates to the back, and may be associated with nausea or vomiting. Severe pancreatitis may lead to bleeding or perforation of the pancreas resulting in shock or a systemic inflammatory response syndrome, bruising of the flanks or the region around the belly button. These severe complications are often managed in an intensive care unit. In pancreatitis, enzymes of the exocrine pancreas damage the structure and tissue of the pancreas. Detection of some of these enzymes, such as amylase and lipase in the blood, along with symptoms and findings on medical imaging such as ultrasound or a CT scan, are often used to indicate that a person has pancreatitis. Pancreatitis is often managed medically with pain relief, and monitoring to prevent or manage shock, and management of any identified underlying causes. This may include removal of gallstones, lowering of blood triglyceride or glucose levels, the use of corticosteroids for autoimmune pancreatitis, and the cessation of any medication triggers. Chronic pancreatitis refers to the development of pancreatitis over time. It shares many similar causes, with the most common being chronic alcohol use, with other causes including recurrent acute episodes and cystic fibrosis. Abdominal pain, characteristically relieved by sitting forward or drinking alcohol, is the most common symptom. When the digestive function of the pancreas is severely affected, this may lead to problems with fat digestion and the development of steatorrhoea; when the endocrine function is affected, this may lead to diabetes. Chronic pancreatitis is investigated in a similar way to acute pancreatitis. In addition to management of pain and nausea, and management of any identified causes (which may include alcohol cessation), because of the digestive role of the pancreas, enzyme replacement may be needed to prevent malabsorption. Cancer Pancreatic cancers, particularly the most common type, pancreatic adenocarcinoma, remain very difficult to treat, and are mostly diagnosed only at a stage that is too late for surgery, which is the only curative treatment. Pancreatic cancer is rare in people younger than 40 and the median age of diagnosis is 71. Risk factors include chronic pancreatitis, older age, smoking, obesity, diabetes, and certain rare genetic conditions including multiple endocrine neoplasia type 1, hereditary nonpolyposis colon cancer and dysplastic nevus syndrome among others. About 25% of cases are attributable to tobacco smoking, while 5–10% of cases are linked to inherited genes. Pancreatic adenocarcinoma is the most common form of pancreatic cancer, and is cancer arising from the exocrine digestive part of the pancreas. Most occur in the head of the pancreas. Symptoms tend to arise late in the course of the cancer, when it causes abdominal pain, weight loss, or yellowing of the skin (jaundice). Jaundice occurs when the outflow of bile is blocked by the cancer. Other less common symptoms include nausea, vomiting, pancreatitis, diabetes or recurrent venous thrombosis. Pancreatic cancer is usually diagnosed by medical imaging in the form of an ultrasound or CT scan with contrast enhancement. An endoscopic ultrasound may be used if a tumour is being considered for surgical removal, and biopsy guided by ERCP or ultrasound can be used to confirm an uncertain diagnosis. Because of the late development of symptoms, most cancer presents at an advanced stage. Only 10 to 15% of tumours are suitable for surgical resection. , when chemotherapy is given the FOLFIRINOX regimen containing fluorouracil, irinotecan, oxaliplatin and leucovorin has been shown to extend survival beyond traditional gemcitabine regimens. For the most part, treatment is palliative, focus on the management of symptoms that develop. This may include management of itch, a choledochojejunostomy or the insertion of stents with ERCP to facilitate the drainage of bile, and medications to help control pain. In the United States pancreatic cancer is the fourth most common cause of deaths due to cancer. The disease occurs more often in the developed world, which had 68% of new cases in 2012. Pancreatic adenocarcinoma typically has poor outcomes with the average percentage alive for at least one and five years after diagnosis being 25% and 5% respectively. In localized disease where the cancer is small (< 2 cm) the number alive at five years is approximately 20%. There are several types of pancreatic cancer, involving both the endocrine and exocrine tissue. The many types of pancreatic endocrine tumors are all uncommon or rare, and have varied outlooks. However the incidence of these cancers has been rising sharply; it is not clear to what extent this reflects increased detection, especially through medical imaging, of tumors that would be very slow to develop. Insulinomas (largely benign) and gastrinomas are the most common types. For those with neuroendocrine cancers the number alive after five years is much better at 65%, varying considerably with type. A solid pseudopapillary tumour is a low-grade malignant tumour of the pancreas of papillary architecture that typically afflicts young women. Diabetes mellitus Type 1 diabetes Diabetes mellitus type 1 is a chronic autoimmune disease in which the immune system attacks the insulin-secreting beta cells of the pancreas. Insulin is needed to keep blood sugar levels within optimal ranges, and its lack can lead to high blood sugar. As an untreated chronic condition, complications including accelerated vascular disease, diabetic retinopathy, kidney disease and neuropathy can result. In addition, if there is not enough insulin for glucose to be used within cells, the medical emergency diabetic ketoacidosis, which is often the first symptom that a person with type 1 diabetes may have, can result. Type 1 diabetes can develop at any age but is most often diagnosed before age 40. For people living with type 1 diabetes, insulin injections are critical for survival. An experimental procedure to treat type 1 diabetes is pancreas transplantation or isolated transplantation of islet cells to supply a person with functioning beta cells. Type 2 diabetes Diabetes mellitus type 2 is the most common form of diabetes. The causes for high blood sugar in this form of diabetes usually are a combination of insulin resistance and impaired insulin secretion, with both genetic and environmental factors playing a role in the development of the disease. Over time, pancreatic beta cells may become "exhausted" and less functional. The management of type 2 diabetes involves a combination of lifestyle measures, medications if required and potentially insulin. With relevance to the pancreas, several medications act to enhance the secretion of insulin from beta cells, particularly sulphonylureas, which act directly on beta cells; incretins which replicate the action of the hormones glucagon-like peptide 1, increasing the secretion of insulin from beta cells after meals, and are more resistant to breakdown; and DPP-4 inhibitors, which slow the breakdown of incretins. Removal It is possible for a person to live without a pancreas, provided that the person takes insulin for proper regulation of blood glucose concentration and pancreatic enzyme supplements to aid digestion. History The pancreas was first identified by Herophilus (335–280 BC), a Greek anatomist and surgeon. A few hundred years later, Rufus of Ephesus, another Greek anatomist, gave the pancreas its name. Etymologically, the term "pancreas", a modern Latin adaptation of Greek πάγκρεας, [πᾶν ("all", "whole"), and κρέας ("flesh")], originally means sweetbread, although literally meaning all-flesh, presumably because of its fleshy consistency. It was only in 1889 when Oskar Minkowski discovered that removing the pancreas from a dog caused it to become diabetic. Insulin was later isolated from pancreatic islets by Frederick Banting and Charles Best in 1921. The way the tissue of the pancreas has been viewed has also changed. Previously, it was viewed using simple staining methods such as H&E stains. Now, immunohistochemistry can be used to more easily differentiate cell types. This involves visible antibodies to the products of certain cell types, and helps identify with greater ease cell types such as alpha and beta cells. Other animals Pancreatic tissue is present in all vertebrates, but its precise form and arrangement varies widely. There may be up to three separate pancreases, two of which arise from the pancreatic bud, and the other dorsally. In most species (including humans), these "fuse" in the adult, but there are several exceptions. Even when a single pancreas is present, two or three pancreatic ducts may persist, each draining separately into the duodenum (or equivalent part of the foregut). Birds, for example, typically have three such ducts. In teleost fish, and a few other species (such as rabbits), there is no discrete pancreas at all, with pancreatic tissue being distributed diffusely across the mesentery and even within other nearby organs, such as the liver or spleen. In a few teleost species, the endocrine tissue has fused to form a distinct gland within the abdominal cavity, but otherwise it is distributed among the exocrine components. The most primitive arrangement, however, appears to be that of lampreys and lungfish, in which pancreatic tissue is found as a number of discrete nodules within the wall of the gut itself, with the exocrine portions being little different from other glandular structures of the intestine. Cuisine The pancreas of calf (ris de veau) or lamb (ris d'agneau), and, less commonly, of beef or pork, are used as food under the culinary name of sweetbread. Additional images	Biology and health sciences	Digestive system	null
38306	https://en.wikipedia.org/wiki/Pegmatite	Pegmatite	A pegmatite is an igneous rock showing a very coarse texture, with large interlocking crystals usually greater in size than and sometimes greater than . Most pegmatites are composed of quartz, feldspar, and mica, having a similar silicic composition to granite. However, rarer intermediate composition and mafic pegmatites are known. Many of the world's largest crystals are found within pegmatites. These include crystals of microcline, quartz, mica, spodumene, beryl, and tourmaline. Some individual crystals are over long. Most pegmatites are thought to form from the last fluid fraction of a large crystallizing magma body. This residual fluid is highly enriched in volatiles and trace elements, and its very low viscosity allows components to migrate rapidly to join an existing crystal rather than coming together to form new crystals. This allows a few very large crystals to form. While most pegmatites have a simple composition of minerals common in ordinary igneous rock, a few pegmatites have a complex composition, with numerous unusual minerals of rare elements. These complex pegmatites are mined for lithium, beryllium, boron, fluorine, tin, tantalum, niobium, rare earth elements, uranium, and other valuable commodities. Etymology The word pegmatite derives from Homeric Greek, πήγνυμι (pēgnymi), which means “to bind together”, in reference to the intertwined crystals of quartz and feldspar in the texture known as graphic granite. The term was first used by René Just Haüy in 1822 as a synonym for graphic granite. Wilhelm Karl Ritter von Haidinger first used the term in its present meaning in 1845. General description Pegmatites are exceptionally coarse-grained igneous rocks composed of interlocking crystals, with individual crystals usually over in size and sometimes exceeding . Most pegmatites have a composition similar to granite, so that their most common minerals are quartz, feldspar, and mica. However, other pegmatite compositions are known, including compositions similar to nepheline syenite or gabbro. The term pegmatite is thus purely a textural description. Geologists typically prefix the term with a compositional description, so that granitic pegmatite is a pegmatite with the composition of granite while nepheline syenite pegmatite is a pegmatite with the composition of nepheline syenite. However, the British Geological Survey (BGS) discourages this usage, preferring terms like biotite-quartz-feldspar pegmatite for a pegmatite with a typical granitic composition, dominated by feldspar with lesser quartz and biotite. Under BGS terminology, a pegmatitic rock (for example, a pegmatitic gabbro) is a coarse-grained rock containing patches of much coarser-grained rock of essentially the same composition. Individual crystals in pegmatites can be enormous in size. It is likely that the largest crystals ever found were feldspar crystals in pegmatites from Karelia with masses of thousands of tons. Quartz crystals with masses measured in thousands of pounds and micas over across and thick have been found. Spodumene crystals over long have been found in the Black Hills of South Dakota, and beryl crystals long and in diameter have been found at Albany, Maine. The largest beryl crystal ever found was from Malakialina on Madagascar, weighing about 380 tons, with a length of and a crosscut of . Pegmatite bodies are usually of minor size compared to typical intrusive rock bodies. Pegmatite body size is on the order of magnitude of one to a few hundred meters. Compared to typical igneous rocks they are rather inhomogeneous and may show zones with different mineral assemblages. Crystal size and mineral assemblages are usually oriented parallel to the wall rock or even concentric for pegmatite lenses. Classification Modern pegmatite classification schemes are strongly influenced by the depth-zone classification of granitic rocks published by Buddington (1959), and the Ginsburg & Rodionov (1960) and Ginsburg et al. (1979) classification which categorized pegmatites according to their depth of emplacement and relationship to metamorphism and granitic plutons. Cerny’s (1991) revision of that classification scheme is widely used, Cerny’s (1991) pegmatite classification, which is a combination of emplacement depth, metamorphic grade and minor element content, has provided significant insight into the origin of pegmatitic melts and their relative degrees of fractionation. Granitic pegmatites are commonly ranked into three hierarchies (class – family – type – subtype) depending upon their mineralogical-geochemical characteristics and depth of emplacement according to Cerny (1991). Classes are Abyssal, Muscovite, Rare-Element and Miarolitic. The Rare-Element Class is subdivided based on composition into LCT and NYF families: LCT for Lithium, Cesium, and Tantalum enrichment and NYF for Niobium, Yttrium, and Fluorine enrichment. Most authors classify pegmatites according to LCT- and NYF-types and subtypes. Another important contribution of the classification is the petrogenetic component of the classification, which shows the association of LCT pegmatites with mainly orogenic plutons, and NYF pegmatites with mainly anorogenic plutons. Lately, there have been a few attempts to create a new classification for pegmatites less dependent on mineralogy and more reflective of their geological setting. On this issue, one of the most notable efforts on this matter is Wise's (2022) pegmatite classification, which focuses mostly on the source of the magma from which the pegmatite crystalizes. Petrology Pegmatites form under conditions in which the rate of new crystal nucleation is much slower than the rate of crystal growth. Large crystals are favored. In normal igneous rocks, coarse texture is a result of slow cooling deep underground. It is not clear if pegmatite forms by slow or rapid cooling. In some studies, crystals in pegmatitic conditions have been recorded to grow at a rate ranging from 1 m to 10 m per day. Pegmatites are the last part of a magma body to crystallize. This final fluid fraction is enriched in volatile and trace elements. The residual magma undergoes phase separation into a melt phase and a hydrous fluid phase saturated with silica, alkalis, and other elements. Such phase separation requires formation from a wet magma, rich enough in water to saturate before more than two-thirds of the magma is crystallized. Otherwise, the separation of the fluid phase is difficult to explain. Granite requires a water content of 4 wt% at a pressure of , but only 1.5 wt% at for phase separation to take place. The volatiles (primarily water, borates, fluorides, chlorides, and phosphates) are concentrated in the hydrous phase, greatly lowering its viscosity. The silica in the hydrous phase is completely depolymerized, existing almost entirely as orthosilicate, with all oxygen bridges between silicon ions broken. The low viscosity promotes rapid diffusion through the fluid, allowing growth of large crystals. When this hydrous fluid is injected into the surrounding country rock, minerals crystallize from the outside in to form a zoned pegmatite, with different minerals predominating in concentric zones. A typical sequence of deposition begins with microcline and quartz, with minor schorl and garnet. This is followed by deposition of albite, lepidolite, gem tourmaline, beryl, spodumene, amblygonite, topaz, apatite, and fluorite, which may partially replace some of the minerals in the earlier zone. The center of the pegmatite may have cavities lined with spectacular gemstone crystals. Some pegmatites have more complex zoning. Five distinct zones are recognized in the Harding Pegmatite in the Picuris Mountains of northern New Mexico, US. These are: A white border rind of fine-grained quartz-albite muscovite pegmatite. A continuous layer of very coarse quartz, albite, and muscovite. This zone also contains microcline, and has abundant accessory apatite, beryl, and tantalite. Beryl is occasionally very coarse and abundant. A continuous layer of massive quartz. This zone is also rich in muscovite, microcline, and cleavelandite. A spectacular quartz and lath-spodumene zone. The spodumene occurs as blade-like crystals, sometimes of enormous size, mostly oriented at random but sometimes arranged to form a comb-like structure. Accessory minerals are beryl, apatite, microcline, and tantalum-niobium minerals, especially in the lower part of this zone. There is some pseudomorphic replacement of spodumene by rose muscovite and quartz by cleavelandite. The core of the pegmatite, known as "spotted rock", which is relatively fine-grained spodumene, microcline, and quartz, with accompanying finer-grained albite, lithium-bearing muscovite, lepidolite, microlite, and tantalite. Much of the spodumene and microcline have been extensively corroded and replaced by fine-grained micas. Large crystals nucleate on the margins of pegmatites, becoming larger as they grow inward. These include very large conical alkali feldspar crystals. Aplites are commonly present. These may cut across the pegmatite, but also form zones or irregular patches around coarser material. The aplites are often layered, showing evidence of deformation. Xenoliths may be found in the body of the pegmatite, but their original mineral content is replaced by quartz and alkali feldspar, so that they are difficult to distinguish from the surrounding pegmatite. Pegmatite also commonly replaces part of the surrounding country rock. Because pegmatites likely crystallize from a fluid-dominated phase, rather than a melt phase, they straddle the boundary between hydrothermal mineral deposits and igneous intrusions. Although there is broad agreement on the basic mechanisms by which they form, the details of pegmatite formation remain enigmatic. Pegmatites have characteristics inconsistent with other igneous intrusions. They are not porphyritic, and show no chilled margin. On the contrary, the largest crystals are often found on the margins of the pegmatite body. While aplites are sometimes found on the margins, they are as likely to occur within the body of the pegmatite. The crystals are never aligned in a way that would indicate flow, but are perpendicular to the walls. This implies formation in a static environment. Some pegmatities take the form of isolated pods, with no obvious feeder conduit. As a result, metamorphic or metasomatic origins have sometimes been suggested for pegmatites. A metamorphic pegmatite would be formed by removal of volatiles from metamorphic rocks, particularly felsic gneiss, to liberate the right constituents and water, at the right temperature. A metasomatic pegmatite would be formed by hydrothermal circulation of hot alteration fluids upon a rock mass, with bulk chemical and textural change. Metasomatism is currently not favored as a mechanism for pegmatite formation and it is likely that metamorphism and magmatism are both contributors toward the conditions necessary for pegmatite genesis. Mineralogy Most pegmatites have a simple composition, often being composed entirely of minerals common in granite, such as feldspar, mica, and quartz. The feldspar and quartz often show graphic texture. Rarely, pegmatites are extremely enriched in incompatible elements, such as lithium, caesium, beryllium, tin, niobium, zirconium, uranium, thorium, boron, phosphorus, and fluorine. These complex pegmatites contain unusual minerals of these elements, such as beryl, spodumene, lepidolite, amblygonite, topaz, apatite, fluorite, tourmaline, triphylite, columbite, monazite, and molybdenite. Some of these can be important ore minerals. Some gemstones, such as emerald, are found almost exclusively in pegmatites. Nepheline syenite pegmatites typically contain zirconium, titanium, and rare earth element minerals. Gabbroic pegmatites typically consist of exceptionally coarse interlocking pyroxene and plagioclase. Geochemistry Pegmatites are enriched in volatile and incompatible elements, consistent with their likely origin as the final melt fraction of a crystallizing body of magma. However, it is difficult to get a representative composition of a pegmatite, due to the large size of the constituent mineral crystals. Hence, pegmatite is often characterised by sampling the individual minerals that compose the pegmatite, and comparisons are made according to mineral chemistry. A common error is to assume that the wall zone is a chilled margin whose composition is representative of the original melt. Pegmatites derived from batholiths can be divided into a family of NYF pegmatites, characterized by progressive enrichment in niobium, yttrium, and fluorine as well as enrichment in beryllium, rare earth elements, scandium, titanium, zirconium, thorium, and uranium; and a family of LCT pegmatites, characterized by progressive accumulation of lithium, caesium, and tantalum, as well as enrichment in rubidium, beryllium, tin, barium, phosphorus, and fluorine. The NYF pegmatites likely fractionated from A- to I-type granites that were relatively low in aluminium (subaluminous to metaluminous granites). These granites originated from depleted crust or mantle rock. LCT pegmatites most likely formed from S-type granites or possibly I-type granites, with a higher aluminium content (peraluminous granites). Intermediate pegmatites (NYF + LCT pegmatites) are known and may have formed by contamination of an initially NYF magma body with melted undepleted supracrustral rock. Economic importance Pegmatites often contain rare elements and gemstones. Examples include aquamarine, tourmaline, topaz, fluorite, apatite, and corundum, often along with tin, rare earth, and tungsten minerals, among others. Pegmatites have been mined for both quartz and feldspar. For quartz mining, pegmatites with central quartz masses have been of particular interest. Pegmatites are the primary source of lithium either as spodumene, lithiophyllite or usually from lepidolite. The primary source for caesium is pollucite, a mineral from a zoned pegmatite. The majority of the world's beryllium is sourced from non-gem quality beryl within pegmatite. Tantalum, niobium, and rare-earth elements are sourced from a few pegmatites worldwide, such as the Greenbushes Pegmatite, the Kibara Belt of Rwanda and Democratic Republic of the Congo, the Kenticha mine of Ethiopia the Alto Ligonha Province of Mozambique, and the Mibra (Volta) mine of Minas Gerais, Brazil. Occurrence Notable pegmatite occurrences are found worldwide within the major cratons, and within greenschist-facies metamorphic belts. However, pegmatite localities are only well recorded when economic mineralisation is found. Pegmatites are found as irregular dikes, sills, or veins, and are most common at the margins of batholiths (great masses of intrusive igneous rock). Most are closely related spatially and genetically to large intrusions. They may take the form of veins or dikes in the intrusion itself, but more commonly, they extend into the surrounding country rock, especially above the intrusion. Some pegmatites surrounded by metamorphic rock have no obvious connection to a larger intrusion. Pegmatites in low-grade metamorphic rock tend to be dominated by quartz and carbonate minerals. Pegmatites in metamorphic rock of higher grade are dominted by alkali feldspar. Gabbroic pegmatites typically occur as lenses within bodies of gabbro or diabase. Nepheline syenite pegmatites are common in alkaline igneous complexes.	Physical sciences	Igneous rocks	Earth science
38310	https://en.wikipedia.org/wiki/Cannabis	Cannabis	Cannabis () is a genus of flowering plants in the family Cannabaceae that is widely accepted as being indigenous to and originating from the continent of Asia. However, the number of species is disputed, with as many as three species being recognized: Cannabis sativa, C. indica, and C. ruderalis. Alternatively, C. ruderalis may be included within C. sativa, or all three may be treated as subspecies of C. sativa, or C. sativa may be accepted as a single undivided species. The plant is also known as hemp, although this term is usually used to refer only to varieties cultivated for non-drug use. Hemp has long been used for fibre, seeds and their oils, leaves for use as vegetables, and juice. Industrial hemp textile products are made from cannabis plants selected to produce an abundance of fibre. Cannabis also has a long history of being used for medicinal purposes, and as a recreational drug known by several slang terms, such as marijuana, pot or weed. Various cannabis strains have been bred, often selectively to produce high or low levels of tetrahydrocannabinol (THC), a cannabinoid and the plant's principal psychoactive constituent. Compounds such as hashish and hash oil are extracted from the plant. More recently, there has been interest in other cannabinoids like cannabidiol (CBD), cannabigerol (CBG), and cannabinol (CBN). Etymology Cannabis is a Scythian word. The ancient Greeks learned of the use of cannabis by observing Scythian funerals, during which cannabis was consumed. In Akkadian, cannabis was known as qunubu (). The word was adopted in to the Hebrew language as qaneh bosem (). Description Cannabis is an annual, dioecious, flowering herb. The leaves are palmately compound or digitate, with serrate leaflets. The first pair of leaves usually have a single leaflet, the number gradually increasing up to a maximum of about thirteen leaflets per leaf (usually seven or nine), depending on variety and growing conditions. At the top of a flowering plant, this number again diminishes to a single leaflet per leaf. The lower leaf pairs usually occur in an opposite leaf arrangement and the upper leaf pairs in an alternate arrangement on the main stem of a mature plant. The leaves have a peculiar and diagnostic venation pattern (which varies slightly among varieties) that allows for easy identification of Cannabis leaves from unrelated species with similar leaves. As is common in serrated leaves, each serration has a central vein extending to its tip, but in Cannabis this originates from lower down the central vein of the leaflet, typically opposite to the position of the second notch down. This means that on its way from the midrib of the leaflet to the point of the serration, the vein serving the tip of the serration passes close by the intervening notch. Sometimes the vein will pass tangentially to the notch, but often will pass by at a small distance; when the latter happens a spur vein (or occasionally two) branches off and joins the leaf margin at the deepest point of the notch. Tiny samples of Cannabis also can be identified with precision by microscopic examination of leaf cells and similar features, requiring special equipment and expertise. Reproduction All known strains of Cannabis are wind-pollinated and the fruit is an achene. Most strains of Cannabis are short day plants, with the possible exception of C. sativa subsp. sativa var. spontanea (= C. ruderalis), which is commonly described as "auto-flowering" and may be day-neutral. Cannabis is predominantly dioecious, having imperfect flowers, with staminate "male" and pistillate "female" flowers occurring on separate plants. "At a very early period the Chinese recognized the Cannabis plant as dioecious", and the (c. 3rd century BCE) Erya dictionary defined xi 枲 "male Cannabis" and fu 莩 (or ju 苴) "female Cannabis". Male flowers are normally borne on loose panicles, and female flowers are borne on racemes. Many monoecious varieties have also been described, in which individual plants bear both male and female flowers. (Although monoecious plants are often referred to as "hermaphrodites", true hermaphrodites – which are less common in Cannabis – bear staminate and pistillate structures together on individual flowers, whereas monoecious plants bear male and female flowers at different locations on the same plant.) Subdioecy (the occurrence of monoecious individuals and dioecious individuals within the same population) is widespread. Many populations have been described as sexually labile. As a result of intensive selection in cultivation, Cannabis exhibits many sexual phenotypes that can be described in terms of the ratio of female to male flowers occurring in the individual, or typical in the cultivar. Dioecious varieties are preferred for drug production, where the fruits (produced by female flowers) are used. Dioecious varieties are also preferred for textile fiber production, whereas monoecious varieties are preferred for pulp and paper production. It has been suggested that the presence of monoecy can be used to differentiate licit crops of monoecious hemp from illicit drug crops, but sativa strains often produce monoecious individuals, which is possibly as a result of inbreeding. Sex determination Cannabis has been described as having one of the most complicated mechanisms of sex determination among the dioecious plants. Many models have been proposed to explain sex determination in Cannabis. Based on studies of sex reversal in hemp, it was first reported by K. Hirata in 1924 that an XY sex-determination system is present. At the time, the XY system was the only known system of sex determination. The X:A system was first described in Drosophila spp in 1925. Soon thereafter, Schaffner disputed Hirata's interpretation, and published results from his own studies of sex reversal in hemp, concluding that an X:A system was in use and that furthermore sex was strongly influenced by environmental conditions. Since then, many different types of sex determination systems have been discovered, particularly in plants. Dioecy is relatively uncommon in the plant kingdom, and a very low percentage of dioecious plant species have been determined to use the XY system. In most cases where the XY system is found it is believed to have evolved recently and independently. Since the 1920s, a number of sex determination models have been proposed for Cannabis. Ainsworth describes sex determination in the genus as using "an X/autosome dosage type". The question of whether heteromorphic sex chromosomes are indeed present is most conveniently answered if such chromosomes were clearly visible in a karyotype. Cannabis was one of the first plant species to be karyotyped; however, this was in a period when karyotype preparation was primitive by modern standards. Heteromorphic sex chromosomes were reported to occur in staminate individuals of dioecious "Kentucky" hemp, but were not found in pistillate individuals of the same variety. Dioecious "Kentucky" hemp was assumed to use an XY mechanism. Heterosomes were not observed in analyzed individuals of monoecious "Kentucky" hemp, nor in an unidentified German cultivar. These varieties were assumed to have sex chromosome composition XX. According to other researchers, no modern karyotype of Cannabis had been published as of 1996. Proponents of the XY system state that Y chromosome is slightly larger than the X, but difficult to differentiate cytologically. More recently, Sakamoto and various co-authors have used random amplification of polymorphic DNA (RAPD) to isolate several genetic marker sequences that they name Male-Associated DNA in Cannabis (MADC), and which they interpret as indirect evidence of a male chromosome. Several other research groups have reported identification of male-associated markers using RAPD and amplified fragment length polymorphism. Ainsworth commented on these findings, stating, Environmental sex determination is known to occur in a variety of species. Many researchers have suggested that sex in Cannabis is determined or strongly influenced by environmental factors. Ainsworth reviews that treatment with auxin and ethylene have feminizing effects, and that treatment with cytokinins and gibberellins have masculinizing effects. It has been reported that sex can be reversed in Cannabis using chemical treatment. A polymerase chain reaction-based method for the detection of female-associated DNA polymorphisms by genotyping has been developed. Chemistry Cannabis plants produce a large number of chemicals as part of their defense against herbivory. One group of these is called cannabinoids, which induce mental and physical effects when consumed. Cannabinoids, terpenes, terpenoids, and other compounds are secreted by glandular trichomes that occur most abundantly on the floral calyxes and bracts of female plants. Genetics Cannabis, like many organisms, is diploid, having a chromosome complement of 2n=20, although polyploid individuals have been artificially produced. The first genome sequence of Cannabis, which is estimated to be 820 Mb in size, was published in 2011 by a team of Canadian scientists. Taxonomy The genus Cannabis was formerly placed in the nettle family (Urticaceae) or mulberry family (Moraceae), and later, along with the genus Humulus (hops), in a separate family, the hemp family (Cannabaceae sensu stricto). Recent phylogenetic studies based on cpDNA restriction site analysis and gene sequencing strongly suggest that the Cannabaceae sensu stricto arose from within the former family Celtidaceae, and that the two families should be merged to form a single monophyletic family, the Cannabaceae sensu lato. Various types of Cannabis have been described, and variously classified as species, subspecies, or varieties: plants cultivated for fiber and seed production, described as low-intoxicant, non-drug, or fiber types. plants cultivated for drug production, described as high-intoxicant or drug types. escaped, hybridised, or wild forms of either of the above types. Cannabis plants produce a unique family of terpeno-phenolic compounds called cannabinoids, some of which produce the "high" which may be experienced from consuming marijuana. There are 483 identifiable chemical constituents known to exist in the cannabis plant, and at least 85 different cannabinoids have been isolated from the plant. The two cannabinoids usually produced in greatest abundance are cannabidiol (CBD) and/or Δ9-tetrahydrocannabinol (THC), but only THC is psychoactive. Since the early 1970s, Cannabis plants have been categorized by their chemical phenotype or "chemotype", based on the overall amount of THC produced, and on the ratio of THC to CBD. Although overall cannabinoid production is influenced by environmental factors, the THC/CBD ratio is genetically determined and remains fixed throughout the life of a plant. Non-drug plants produce relatively low levels of THC and high levels of CBD, while drug plants produce high levels of THC and low levels of CBD. When plants of these two chemotypes cross-pollinate, the plants in the first filial (F1) generation have an intermediate chemotype and produce intermediate amounts of CBD and THC. Female plants of this chemotype may produce enough THC to be utilized for drug production. Whether the drug and non-drug, cultivated and wild types of Cannabis constitute a single, highly variable species, or the genus is polytypic with more than one species, has been a subject of debate for well over two centuries. This is a contentious issue because there is no universally accepted definition of a species. One widely applied criterion for species recognition is that species are "groups of actually or potentially interbreeding natural populations which are reproductively isolated from other such groups." Populations that are physiologically capable of interbreeding, but morphologically or genetically divergent and isolated by geography or ecology, are sometimes considered to be separate species. Physiological barriers to reproduction are not known to occur within Cannabis, and plants from widely divergent sources are interfertile. However, physical barriers to gene exchange (such as the Himalayan mountain range) might have enabled Cannabis gene pools to diverge before the onset of human intervention, resulting in speciation. It remains controversial whether sufficient morphological and genetic divergence occurs within the genus as a result of geographical or ecological isolation to justify recognition of more than one species. Early classifications The genus Cannabis was first classified using the "modern" system of taxonomic nomenclature by Carl Linnaeus in 1753, who devised the system still in use for the naming of species. He considered the genus to be monotypic, having just a single species that he named Cannabis sativa L. Linnaeus was familiar with European hemp, which was widely cultivated at the time. This classification was supported by Christiaan Hendrik Persoon (in 1807), Lindley (in 1838) and De Candollee (in 1867). These first classification attempts resulted in a four group division: Kif (southern hemp - psychoactive) Vulgaris (intermediate - psychoactive and fiber) Pedemontana (northern hemp - fiber) Chinensis (northern hemp - fiber) In 1785, evolutionary biologist Jean-Baptiste de Lamarck published a description of a second species of Cannabis, which he named Cannabis indica Lam. Lamarck based his description of the newly named species on morphological aspects (trichomes, leaf shape) and geographic localization of plant specimens collected in India. He described C. indica as having poorer fiber quality than C. sativa, but greater utility as an inebriant. Also, C. indica was considered smaller, by Lamarck. Also, woodier stems, alternate ramifications of the branches, narrow leaflets, and a villous calyx in the female flowers were characteristics noted by the botanist. In 1843, William O’Shaughnessy, used "Indian hemp (C. indica)" in a work title. The author claimed that this choice wasn't based on a clear distinction between C. sativa and C. indica, but may have been influenced by the choice to use the term "Indian hemp" (linked to the plant's history in India), hence naming the species as indica. Additional Cannabis species were proposed in the 19th century, including strains from China and Vietnam (Indo-China) assigned the names Cannabis chinensis Delile, and Cannabis gigantea Delile ex Vilmorin. However, many taxonomists found these putative species difficult to distinguish. In the early 20th century, the single-species concept (monotypic classification) was still widely accepted, except in the Soviet Union, where Cannabis continued to be the subject of active taxonomic study. The name Cannabis indica was listed in various Pharmacopoeias, and was widely used to designate Cannabis suitable for the manufacture of medicinal preparations. 20th century In 1924, Russian botanist D.E. Janichevsky concluded that ruderal Cannabis in central Russia is either a variety of C. sativa or a separate species, and proposed C. sativa L. var. ruderalis Janisch, and Cannabis ruderalis Janisch, as alternative names. In 1929, renowned plant explorer Nikolai Vavilov assigned wild or feral populations of Cannabis in Afghanistan to C. indica Lam. var. kafiristanica Vav., and ruderal populations in Europe to C. sativa L. var. spontanea Vav. Vavilov, in 1931, proposed a three species system, independently reinforced by Schultes et al (1975) and Emboden (1974): C. sativa, C. indica and C. ruderalis. In 1940, Russian botanists Serebriakova and Sizov proposed a complex poly-species classification in which they also recognized C. sativa and C. indica as separate species. Within C. sativa they recognized two subspecies: C. sativa L. subsp. culta Serebr. (consisting of cultivated plants), and C. sativa L. subsp. spontanea (Vav.) Serebr. (consisting of wild or feral plants). Serebriakova and Sizov split the two C. sativa subspecies into 13 varieties, including four distinct groups within subspecies culta. However, they did not divide C. indica into subspecies or varieties. Zhukovski, in 1950, also proposed a two-species system, but with C. sativa L. and C. ruderalis. In the 1970s, the taxonomic classification of Cannabis took on added significance in North America. Laws prohibiting Cannabis in the United States and Canada specifically named products of C. sativa as prohibited materials. Enterprising attorneys for the defense in a few drug busts argued that the seized Cannabis material may not have been C. sativa, and was therefore not prohibited by law. Attorneys on both sides recruited botanists to provide expert testimony. Among those testifying for the prosecution was Dr. Ernest Small, while Dr. Richard E. Schultes and others testified for the defense. The botanists engaged in heated debate (outside of court), and both camps impugned the other's integrity. The defense attorneys were not often successful in winning their case, because the intent of the law was clear. In 1976, Canadian botanist Ernest Small and American taxonomist Arthur Cronquist published a taxonomic revision that recognizes a single species of Cannabis with two subspecies (hemp or drug; based on THC and CBD levels) and two varieties in each (domesticated or wild). The framework is thus: C. sativa L. subsp. sativa, presumably selected for traits that enhance fiber or seed production. C. sativa L. subsp. sativa var. sativa, domesticated variety. C. sativa L. subsp. sativa var. spontanea Vav., wild or escaped variety. C. sativa L. subsp. indica (Lam.) Small & Cronq., primarily selected for drug production. C. sativa L. subsp. indica var. indica, domesticated variety. C. sativa subsp. indica var. kafiristanica (Vav.) Small & Cronq, wild or escaped variety. This classification was based on several factors including interfertility, chromosome uniformity, chemotype, and numerical analysis of phenotypic characters. Professors William Emboden, Loran Anderson, and Harvard botanist Richard E. Schultes and coworkers also conducted taxonomic studies of Cannabis in the 1970s, and concluded that stable morphological differences exist that support recognition of at least three species, C. sativa, C. indica, and C. ruderalis. For Schultes, this was a reversal of his previous interpretation that Cannabis is monotypic, with only a single species. According to Schultes' and Anderson's descriptions, C. sativa is tall and laxly branched with relatively narrow leaflets, C. indica is shorter, conical in shape, and has relatively wide leaflets, and C. ruderalis is short, branchless, and grows wild in Central Asia. This taxonomic interpretation was embraced by Cannabis aficionados who commonly distinguish narrow-leafed "sativa" strains from wide-leafed "indica" strains. McPartland's review finds the Schultes taxonomy inconsistent with prior work (protologs) and partly responsible for the popular usage. Continuing research Molecular analytical techniques developed in the late 20th century are being applied to questions of taxonomic classification. This has resulted in many reclassifications based on evolutionary systematics. Several studies of random amplified polymorphic DNA (RAPD) and other types of genetic markers have been conducted on drug and fiber strains of Cannabis, primarily for plant breeding and forensic purposes. Dutch Cannabis researcher E.P.M. de Meijer and coworkers described some of their RAPD studies as showing an "extremely high" degree of genetic polymorphism between and within populations, suggesting a high degree of potential variation for selection, even in heavily selected hemp cultivars. They also commented that these analyses confirm the continuity of the Cannabis gene pool throughout the studied accessions, and provide further confirmation that the genus consists of a single species, although theirs was not a systematic study per se. An investigation of genetic, morphological, and chemotaxonomic variation among 157 Cannabis accessions of known geographic origin, including fiber, drug, and feral populations showed cannabinoid variation in Cannabis germplasm. The patterns of cannabinoid variation support recognition of C. sativa and C. indica as separate species, but not C. ruderalis. C. sativa contains fiber and seed landraces, and feral populations, derived from Europe, Central Asia, and Turkey. Narrow-leaflet and wide-leaflet drug accessions, southern and eastern Asian hemp accessions, and feral Himalayan populations were assigned to C. indica. In 2005, a genetic analysis of the same set of accessions led to a three-species classification, recognizing C. sativa, C. indica, and (tentatively) C. ruderalis. Another paper in the series on chemotaxonomic variation in the terpenoid content of the essential oil of Cannabis revealed that several wide-leaflet drug strains in the collection had relatively high levels of certain sesquiterpene alcohols, including guaiol and isomers of eudesmol, that set them apart from the other putative taxa. A 2020 analysis of single-nucleotide polymorphisms reports five clusters of cannabis, roughly corresponding to hemps (including folk "Ruderalis") folk "Indica" and folk "Sativa". Despite advanced analytical techniques, much of the cannabis used recreationally is inaccurately classified. One laboratory at the University of British Columbia found that Jamaican Lamb's Bread, claimed to be 100% sativa, was in fact almost 100% indica (the opposite strain). Legalization of cannabis in Canada () may help spur private-sector research, especially in terms of diversification of strains. It should also improve classification accuracy for cannabis used recreationally. Legalization coupled with Canadian government (Health Canada) oversight of production and labelling will likely result in more—and more accurate—testing to determine exact strains and content. Furthermore, the rise of craft cannabis growers in Canada should ensure quality, experimentation/research, and diversification of strains among private-sector producers. Popular usage The scientific debate regarding taxonomy has had little effect on the terminology in widespread use among cultivators and users of drug-type Cannabis. Cannabis aficionados recognize three distinct types based on such factors as morphology, native range, aroma, and subjective psychoactive characteristics. "Sativa" is the most widespread variety, which is usually tall, laxly branched, and found in warm lowland regions. "Indica" designates shorter, bushier plants adapted to cooler climates and highland environments. "Ruderalis" is the informal name for the short plants that grow wild in Europe and Central Asia. Mapping the morphological concepts to scientific names in the Small 1976 framework, "Sativa" generally refers to C. sativa subsp. indica var. indica, "Indica" generally refers to C. sativa subsp. i. kafiristanica (also known as afghanica), and "Ruderalis", being lower in THC, is the one that can fall into C. sativa subsp. sativa. The three names fit in Schultes's framework better, if one overlooks its inconsistencies with prior work. Definitions of the three terms using factors other than morphology produces different, often conflicting results. Breeders, seed companies, and cultivators of drug type Cannabis often describe the ancestry or gross phenotypic characteristics of cultivars by categorizing them as "pure indica", "mostly indica", "indica/sativa", "mostly sativa", or "pure sativa". These categories are highly arbitrary, however: one "AK-47" hybrid strain has received both "Best Sativa" and "Best Indica" awards. Phylogeny Cannabis likely split from its closest relative, Humulus (hops), during the mid Oligocene, around 27.8 million years ago according to molecular clock estimates. The centre of origin of Cannabis is likely in the northeastern Tibetan Plateau. The pollen of Humulus and Cannabis are very similar and difficult to distinguish. The oldest pollen thought to be from Cannabis is from Ningxia, China, on the boundary between the Tibetan Plateau and the Loess Plateau, dating to the early Miocene, around 19.6 million years ago. Cannabis was widely distributed over Asia by the Late Pleistocene. The oldest known Cannabis in South Asia dates to around 32,000 years ago. Uses Cannabis is used for a wide variety of purposes. History According to genetic and archaeological evidence, cannabis was first domesticated about 12,000 years ago in East Asia during the early Neolithic period. The use of cannabis as a mind-altering drug has been documented by archaeological finds in prehistoric societies in Eurasia and Africa. The oldest written record of cannabis usage is the Greek historian Herodotus's reference to the central Eurasian Scythians taking cannabis steam baths. His () Histories records, "The Scythians, as I said, take some of this hemp-seed [presumably, flowers], and, creeping under the felt coverings, throw it upon the red-hot stones; immediately it smokes, and gives out such a vapour as no Greek vapour-bath can exceed; the Scyths, delighted, shout for joy." Classical Greeks and Romans also used cannabis. In China, the psychoactive properties of cannabis are described in the Shennong Bencaojing (3rd century AD). Cannabis smoke was inhaled by Daoists, who burned it in incense burners. In the Middle East, use spread throughout the Islamic empire to North Africa. In 1545, cannabis spread to the western hemisphere where Spaniards imported it to Chile for its use as fiber. In North America, cannabis, in the form of hemp, was grown for use in rope, cloth and paper. Cannabinol (CBN) was the first compound to be isolated from cannabis extract in the late 1800s. Its structure and chemical synthesis were achieved by 1940, followed by some of the first preclinical research studies to determine the effects of individual cannabis-derived compounds in vivo. Globally, in 2013, 60,400 kilograms of cannabis were produced legally. Recreational use Cannabis is a popular recreational drug around the world, only behind alcohol, caffeine, and tobacco. In the U.S. alone, it is believed that over 100 million Americans have tried cannabis, with 25 million Americans having used it within the past year. As a drug it usually comes in the form of dried marijuana, hashish, or various extracts collectively known as hashish oil. Normal cognition is restored after approximately three hours for larger doses via a smoking pipe, bong or vaporizer. However, if a large amount is taken orally the effects may last much longer. After 24 hours to a few days, minuscule psychoactive effects may be felt, depending on dosage, frequency and tolerance to the drug. Cannabidiol (CBD), which has no intoxicating effects by itself (although sometimes showing a small stimulant effect, similar to caffeine), is thought to attenuate (i.e., reduce) the anxiety-inducing effects of high doses of THC, particularly if administered orally prior to THC exposure. According to Delphic analysis by British researchers in 2007, cannabis has a lower risk factor for dependence compared to both nicotine and alcohol. However, everyday use of cannabis may be correlated with psychological withdrawal symptoms, such as irritability or insomnia, and susceptibility to a panic attack may increase as levels of THC metabolites rise. Cannabis withdrawal symptoms are typically mild and are not life-threatening. Risk of adverse outcomes from cannabis use may be reduced by implementation of evidence-based education and intervention tools communicated to the public with practical regulation measures. In 2014 there were an estimated 182.5 million cannabis users worldwide (3.8% of the global population aged 15–64). This percentage did not change significantly between 1998 and 2014. Medical use Medical cannabis (or medical marijuana) refers to the use of cannabis and its constituent cannabinoids, in an effort to treat disease or improve symptoms. Cannabis is used to reduce nausea and vomiting during chemotherapy, to improve appetite in people with HIV/AIDS, and to treat chronic pain and muscle spasms. Cannabinoids are under preliminary research for their potential to affect stroke. Evidence is lacking for depression, anxiety, attention deficit hyperactivity disorder, Tourette syndrome, post-traumatic stress disorder, and psychosis. Two extracts of cannabis – dronabinol and nabilone – are approved by the FDA as medications in pill form for treating the side effects of chemotherapy and AIDS. Short-term use increases both minor and major adverse effects. Common side effects include dizziness, feeling tired, vomiting, and hallucinations. Long-term effects of cannabis are not clear. Concerns including memory and cognition problems, risk of addiction, schizophrenia in young people, and the risk of children taking it by accident. Industrial use (hemp) The term hemp is used to name the durable soft fiber from the Cannabis plant stem (stalk). Cannabis sativa cultivars are used for fibers due to their long stems; Sativa varieties may grow more than six metres tall. However, hemp can refer to any industrial or foodstuff product that is not intended for use as a drug. Many countries regulate limits for psychoactive compound (THC) concentrations in products labeled as hemp. Cannabis for industrial uses is valuable in tens of thousands of commercial products, especially as fibre ranging from paper, cordage, construction material and textiles in general, to clothing. Hemp is stronger and longer-lasting than cotton. It also is a useful source of foodstuffs (hemp milk, hemp seed, hemp oil) and biofuels. Hemp has been used by many civilizations, from China to Europe (and later North America) during the last 12,000 years. In modern times novel applications and improvements have been explored with modest commercial success. In the US, "industrial hemp" is classified by the federal government as cannabis containing no more than 0.3% THC by dry weight. This classification was established in the 2018 Farm Bill and was refined to include hemp-sourced extracts, cannabinoids, and derivatives in the definition of hemp. Ancient and religious uses The Cannabis plant has a history of medicinal use dating back thousands of years across many cultures. The Yanghai Tombs, a vast ancient cemetery (54 000 m2) situated in the Turfan district of the Xinjiang Uyghur Autonomous Region in northwest China, have revealed the 2700-year-old grave of a shaman. He is thought to have belonged to the Jushi culture recorded in the area centuries later in the Hanshu, Chap 96B. Near the head and foot of the shaman was a large leather basket and wooden bowl filled with 789g of cannabis, superbly preserved by climatic and burial conditions. An international team demonstrated that this material contained THC. The cannabis was presumably employed by this culture as a medicinal or psychoactive agent, or an aid to divination. This is the oldest documentation of cannabis as a pharmacologically active agent. The earliest evidence of cannabis smoking has been found in the 2,500-year-old tombs of Jirzankal Cemetery in the Pamir Mountains in Western China, where cannabis residue were found in burners with charred pebbles possibly used during funeral rituals. Settlements which date from c. 2200–1700 BCE in the Bactria and Margiana contained elaborate ritual structures with rooms containing everything needed for making drinks containing extracts from poppy (opium), hemp (cannabis), and ephedra (which contains ephedrine). Although there is no evidence of ephedra being used by steppe tribes, they engaged in cultic use of hemp. Cultic use ranged from Romania to the Yenisei River and had begun by 3rd millennium BC Smoking hemp has been found at Pazyryk. Cannabis is first referred to in Hindu Vedas between 2000 and 1400 BCE, in the Atharvaveda. By the 10th century CE, it has been suggested that it was referred to by some in India as "food of the gods". Cannabis use eventually became a ritual part of the Hindu festival of Holi. One of the earliest to use this plant in medical purposes was Korakkar, one of the 18 Siddhas. The plant is called Korakkar Mooli in the Tamil language, meaning Korakkar's herb. In Buddhism, cannabis is generally regarded as an intoxicant and may be a hindrance to development of meditation and clear awareness. In ancient Germanic culture, Cannabis was associated with the Norse love goddess, Freya. An anointing oil mentioned in Exodus is, by some translators, said to contain Cannabis. In modern times, the Rastafari movement has embraced Cannabis as a sacrament. Elders of the Ethiopian Zion Coptic Church, a religious movement founded in the U.S. in 1975 with no ties to either Ethiopia or the Coptic Church, consider Cannabis to be the Eucharist, claiming it as an oral tradition from Ethiopia dating back to the time of Christ. Like the Rastafari, some modern Gnostic Christian sects have asserted that Cannabis is the Tree of Life. Other organized religions founded in the 20th century that treat Cannabis as a sacrament are the THC Ministry, Cantheism, the Cannabis Assembly and the Church of Cognizance. Since the 13th century CE, cannabis has been used among Sufis – the mystical interpretation of Islam that exerts strong influence over local Muslim practices in Bangladesh, India, Indonesia, Turkey, and Pakistan. Cannabis preparations are frequently used at Sufi festivals in those countries. Pakistan's Shrine of Lal Shahbaz Qalandar in Sindh province is particularly renowned for the widespread use of cannabis at the shrine's celebrations, especially its annual Urs festival and Thursday evening dhamaal sessions – or meditative dancing sessions.	Biology and health sciences	Rosales	null
38323	https://en.wikipedia.org/wiki/Cigar	Cigar	A cigar is a rolled bundle of dried and fermented tobacco leaves made to be smoked. Cigars are produced in a variety of sizes and shapes. Since the 20th century, almost all cigars are made of three distinct components: the filler, the binder leaf which holds the filler together, and a wrapper leaf, which is often the highest quality leaf used. Often there will be a cigar band printed with the cigar manufacturer's logo. Modern cigars can come with two or more bands, especially Cuban cigars, showing Limited Edition (Edición Limitada) bands displaying the year of production. Cigar tobacco is grown in significant quantities primarily in Brazil, Central America (Costa Rica, Ecuador, Guatemala, Honduras, Mexico, Nicaragua, and Panama), and the islands of the Caribbean (Cuba, the Dominican Republic, Haiti, and Puerto Rico); it is also produced in the Eastern United States (mostly in Florida, Kentucky, Tennessee, and Virginia) and in the Mediterranean countries of Italy, Greece, Spain (in the Canary Islands), and Turkey, and to a lesser degree in Indonesia and the Philippines of Southeast Asia. Cigar smoking carries serious health risks, including increased risk of developing various types and subtypes of cancers, respiratory diseases, cardiovascular diseases, cerebrovascular diseases, periodontal diseases, teeth decay and loss, and malignant diseases. In the United States, the tobacco industry and cigar brands have aggressively targeted African Americans and Non-Hispanic Whites with customized advertising techniques and tobacco-related lifestyle magazines since the 1990s. Etymology The word cigar originally derives from the Mayan sikar ("to smoke rolled tobacco leaves"—from si'c, "tobacco"). The Spanish word, "cigarro" spans the gap between the Mayan and modern use. The English word came into general use in 1730. History Although the origins of cigar smoking are unknown, cigar smoking was first observed by European explorers when encountering the indigenous Taino people of Cuba in 1492. While tobacco was widely diffused among many of the Indigenous peoples of the islands of the Caribbean, it was completely unfamiliar to Europeans before the discovery of the New World in the 15th century. The Spanish historian, landowner, and Dominican friar Bartolomé de las Casas vividly described how the first scouts sent by Christopher Columbus into the interior of Cuba found Following the arrival of Europeans with the first wave of European colonization, tobacco became one of the primary products fueling European colonialism, and also became a driving factor in the incorporation of African slave labor. The Spanish introduced tobacco to Europeans in about 1528, and by 1533, Diego Columbus mentioned a tobacco merchant of Lisbon in his will, showing how quickly the traffic had sprung up. The French, Spanish, and Portuguese initially referred to the plant as the "sacred herb" because of its alleged medicinal properties. In time, Spanish and other European sailors adopted the practice of smoking rolls of leaves, as did the Spanish and Portuguese conquistadors. Smoking primitive cigars spread to Spain, Portugal, and eventually France, most probably through Jean Nicot, the French ambassador to Portugal, who gave his name to nicotine. Later, tobacco use spread to the Italian kingdoms, the Dutch Empire, and, after Sir Walter Raleigh's voyages to the Americas, to Great Britain. Tobacco smoking became familiar throughout Europe—in pipes in Britain—by the mid-16th century. Spanish cultivation of tobacco began in earnest in 1531 on the islands of Hispaniola and Santo Domingo. In 1542, tobacco started to be grown commercially in North America, when Spaniards established the first cigar factory in Cuba. Tobacco was originally thought to have medicinal qualities, but some considered it evil. It was denounced by Philip II of Spain and James I of England. Around 1592, the Spanish galleon San Clemente brought of tobacco seed to the Philippines over the Acapulco-Manila trade route. It was distributed among Roman Catholic missionaries, who found excellent climates and soils for growing high-quality tobacco there. The use of the cigar did not become popular until the mid 18th century, and although there are few drawings from this era, there are some reports. It is believed that Israel Putnam brought back a cache of Havana cigars during the Seven Years' War, making cigar smoking popular in the US after the American Revolution. He also brought Cuban tobacco seeds, which he planted in the Hartford area of New England. This reportedly resulted in the development of the renowned shade-grown Connecticut wrapper. Towards the end of the 18th century and in the 19th century, cigar smoking was common, while cigarettes were comparatively rare. Towards the end of the 19th century, Rudyard Kipling wrote his famous smoking poem, The Betrothed (1886). The cigar business was an important industry and factories employed many people before mechanized manufacturing of cigars became practical. Cigar workers in both Cuba and the US were active in labor strikes and disputes from early in the 19th century, and the rise of modern labor unions can be traced to the CMIU and other cigar worker unions. In 1869, Spanish cigar manufacturer Vicente Martinez Ybor moved his Principe de Gales (Prince of Wales) operations from the cigar manufacturing center of Havana, Cuba to Key West, Florida to escape the turmoil of the Ten Years' War. Other manufacturers followed, and Key West became an important cigar manufacturing center. In 1885, Ybor moved again, buying land near the small city of Tampa, Florida and building the largest cigar factory in the world at the time in the new company town of Ybor City. Friendly rival and Flor de Sánchez y Haya owner Ignacio Haya built his factory nearby the same year, and many other cigar manufacturers followed, especially after an 1886 fire that gutted much of Key West. Thousands of Cuban and Spanish tabaqueros came to the area from Key West, Cuba and New York to produce hundreds of millions of cigars annually. Local output peaked in 1929, when workers in Ybor City and West Tampa rolled over 500 million "clear Havana" cigars, earning the town the nickname "Cigar Capital of the World". At its peak, there were 150 cigar factories in Ybor city, but by early in the next decade, nearly all of the factories had closed. Only one company still makes cigars in the Ybor City area, the J. C. Newman Cigar Company, which moved to Tampa from Ohio in 1954 and took over the previous Regensburg cigar factory. The company was continuing to utilize some antique, hand-operated ARENCO and American Machine and Foundry cigarmaking machines from the 1930's. In New York, cigars were made by rollers working in their homes. It was reported that as of 1883, cigars were being manufactured in 127 apartment houses in New York, employing 1,962 families and 7,924 individuals. A state statute banning the practice, passed late that year at the urging of trade unions on the basis that the practice suppressed wages, was ruled unconstitutional less than four months later. The industry, which had relocated to Brooklyn (then a separate municipality) and other places on Long Island while the law was in effect, then returned to New York. As of 1905, there were 80,000 cigar-making operations in the US, most of them small, family-operated shops where cigars were rolled and sold immediately. While most cigars are now made by machine, some, as a matter of prestige and quality, are rolled by hand—especially in Central America and Cuba, as well as in small chinchales in sizable cities in the US. Manufacture Tobacco leaves are harvested and aged using a curing process that combines heat and shade to reduce sugar and water content without causing the larger leaves to rot. This takes between 25 and 45 days, depending upon climatic conditions and the nature of sheds used to store harvested tobacco. Curing varies by type of tobacco and desired leaf color. A slow fermentation follows, where temperature and humidity are controlled to enhance flavor, aroma, and burning characteristics while forestalling rot or disintegration. The leaf will continue to be baled, inspected, un-baled, re-inspected, and baled again during the aging cycle. When it has matured to manufacturer's specifications it is sorted for appearance and overall quality, and used as filler or wrapper accordingly. During this process, leaves are continually moistened to prevent damage. Quality cigars are still handmade. An experienced cigar-roller can produce hundreds of good, nearly identical cigars per day. The rollers keep the tobacco moist—especially the wrapper—and use specially designed crescent-shaped knives, called chavetas, to form the filler and wrapper leaves quickly and accurately. Once rolled, the cigars are stored in wooden forms as they dry, in which their uncapped ends are cut to a uniform size. From this stage, the cigar is a complete product that can be "laid down" and aged for decades if kept as close to 21 °C (70 °F) and 70% relative humidity as possible. Once purchased, proper storage is typically in a specialized cedar-lined wooden humidor. Some cigars, especially premium brands, use different varieties of tobacco for the filler and the wrapper. Long filler cigars are a far higher quality of cigar, using long leaves throughout. These cigars also use a third variety of tobacco leaf, called a "binder", between the filler and the outer wrapper. This permits the makers to use more delicate and attractive leaves as a wrapper. These high-quality cigars almost always blend varieties of tobacco. Even Cuban long-filler cigars will combine tobaccos from different parts of the island to incorporate several different flavors. In low-grade and machine-made cigars, chopped tobacco leaves are used for the filler, and long leaves or a type of "paper" made from reconstituted tobacco pulp is used for the wrapper. Chopped leaves and a pulp wrapper alter the flavor and burning characteristics of the result vis-a-vis handmade cigars. Historically, a lector or reader was employed to entertain cigar factory workers. This practice became obsolete once audiobooks for portable music players became available, but it is still practiced in some Cuban factories. Dominant manufacturers Two firms dominate the cigar industry, Altadis and the Scandinavian Tobacco Group. Altadis, a Spanish-owned private concern, produces cigars in the US, the Dominican Republic, and Honduras, and owns a 50% stake in Corporación Habanos S.A., the state owned national Cuban tobacco company. It also makes cigarettes. The Scandinavian Tobacco Group produces cigars in the Dominican Republic, Honduras, Nicaragua, Indonesia, the Netherlands, Belgium, Denmark and the United States; it also makes pipe tobacco and fine cut tobacco. The Group includes General Cigar Co. The town of Tamboril in Santiago, Dominican Republic is considered by many as today's "Cigar Capital of the World" housing more cigar factories and rollers than anywhere else in the world. According to Cigar Aficionado magazine, 44% of the world's most traded cigars come from the Dominican Republic, the world's largest producer of cigars, especially from the fertile lands of the Cibao capital, where 90% of the factories are located. The area has also been the largest supplier of cigars to the US in the last decades. Families in the cigar industry Nearly all modern premium cigar makers are members of long-established cigar families, or purport to be, most originally rooted in the historic Cuban cigar industry. The art and skill of hand-making premium cigars has been passed from generation to generation. Families are often shown in many cigar advertisements and packaging. In 1992, Cigar Aficionado magazine created the "Cigar Hall of Fame" and recognized the following six individuals: Edgar M. Cullman, Chairman, General Cigar Company, New York, United States Zino Davidoff, Founder, Davidoff et Cie., Geneva, Switzerland Carlos Fuente Sr., Chairman, Tabacalera A. Fuente y Cia., Santiago de los Caballeros, Dominican Republic Frank Llaneza, Chairman, Villazon & Co., Tampa, Florida, United States Stanford J. Newman, Chairman, J.C. Newman Cigar Company, Tampa, Florida, United States Ángel Oliva Sr. (founder); Oliva Tobacco Co., Tampa, Florida, United States Other families in the cigar industry (2015) Manuel Quesada (MATASA Current CEO) Fonseca, Casa Magna, Quesada cigars, Dominican Republic Don José "Pepín" Garcia, Chairman, El Rey de Los Habanos, Miami, Florida, United States Aray Family – Daniel Aray Jr, Grandson of Founder (1952) Jose Aray, ACC Cigars, Guayaquil Ecuador, San Francisco, CA, Miami Florida, Macau SAR, Shanghai China. EPC – Ernesto Perez-Carillo, Founder EPC Cigar Company (2009), Miami, Florida, United States Nestor Miranda – Founder, Miami Cigar Company (1989) Miami, FL, United States Blanco family – Jose "Jochy" Blanco, son of Founder (1936) Jose Arnaldo Blanco Polanco, Tabacalera La Palma, Santiago, Dominican Republic Hermann Dietrich Upmann, founder of the H. Upmann brand 1844 in Cuba Marketing and distribution Pure tobacco, hand rolled cigars are marketed via advertisements, product placement in movies and other media, sporting events, cigar-friendly magazines such as Cigar Aficionado, and cigar dinners. Since handmade cigars are a premium product with a hefty price, advertisements often include depictions of affluence, sensual imagery, and explicit or implied celebrity endorsement. Cigar Aficionado, launched in 1992, presents cigars as symbols of a successful lifestyle, and is a major conduit of advertisements that do not conform to the tobacco industry's voluntary advertisement restrictions since 1965, such as a restriction not to associate smoking with glamour. The magazine also presents pro-smoking arguments at length, and argues that cigars are safer than cigarettes, since they do not have the thousands of chemical additives that cigarette manufactures add to the cutting floor scraps of tobacco used as cigarette filler. The publication also presents arguments that risks are a part of daily life and that (contrary to the evidence discussed in Health effects) cigar smoking has health benefits, that moderation eliminates most or all health risk, and that cigar smokers live to old age, that health research is flawed, and that several health-research results support claims of safety. Like its competitor Smoke, Cigar Aficionado differs from marketing vehicles used for other tobacco products in that it makes cigars the main (but not sole) focus of the magazine, creating a symbiosis between product and lifestyle. In the US, cigars have historically been exempt from many of the marketing regulations that govern cigarettes. For example, the Public Health Cigarette Smoking Act of 1970 exempted cigars from its advertising ban, and cigar ads, unlike cigarette ads, need not mention health risks. As of 2007, cigars were taxed far less than cigarettes, so much so that in many US states, a pack of little cigars cost less than half as much as a pack of cigarettes. It is illegal for minors to purchase cigars and other tobacco products in the US, but laws are unevenly enforced: a 2000 study found that three-quarters of web cigar sites allowed minors to purchase them. In 2009, the US Family Smoking Prevention and Tobacco Control Act provided the Food and Drug Administration regulatory authority over the manufacturing, distribution, and marketing of cigarettes, roll-your-own tobacco and smokeless tobacco. In 2016, a deeming rule extended the FDA's authority to additional tobacco products including cigars, e-cigarettes and hookah. The objective of the law is to reduce the impact of tobacco on public health by preventing Americans from starting to use tobacco products, encourage current users to quit, and decrease the harms of tobacco product use. In the US, inexpensive cigars are sold in convenience stores, gas stations, grocery stores, and pharmacies. Premium cigars are sold in tobacconists, cigar bars, and other specialized establishments. Some cigar stores are part of chains, which have varied in size: in the US, United Cigar Stores was one of only three outstanding examples of national chains in the early 1920s, the others being A&P and Woolworth's. Non-traditional outlets for cigars include hotel shops, restaurants, vending machines and the Internet. Composition Cigars are composed of three types of tobacco leaves, whose variations determine smoking and flavor characteristics: Wrapper A cigar's outermost layer, or wrapper (Spanish: ), is the most expensive component of a cigar. The wrapper determines much of the cigar's character and flavor, and as such its color is often used to describe the cigar as a whole. Wrappers are frequently grown underneath huge canopies made of gauze so as to diffuse direct sunlight and are fermented separately from other rougher cigar components, with a view to the production of a thinly-veined, smooth, supple leaf. Wrapper tobacco produced without the gauze canopies under which "shade grown" leaf is grown, generally more coarse in texture and stronger in flavor, is commonly known as "sun grown". A number of different countries are used for the production of wrapper tobacco, including Cuba, Ecuador, Indonesia, Honduras, Nicaragua, Costa Rica, Brazil, Mexico, Cameroon, and the United States. While dozens of minor wrapper shades have been touted by manufacturers, the seven most common classifications are as follows, ranging from lightest to darkest: Some manufacturers use an alternate designation: In general, dark wrappers add a touch of sweetness, while light ones add a hint of dryness to the taste. Binder Beneath the wrapper is a small bunch of "filler" leaves bound together inside of a leaf called a "binder" (Spanish: ). The binder leaf is typically the sun-saturated leaf from the top part of a tobacco plant and is selected for its elasticity and durability in the rolling process. Unlike the wrapper leaf, which must be uniform in appearance and smooth in texture, the binder leaf may show evidence of physical blemishes or lack uniform coloration. The binder leaf is generally considerably thicker and hardier than the wrapper leaf surrounding it. Filler The bulk of a cigar is "filler"—a bound bunch of tobacco leaves. These leaves are folded by hand to allow air passageways down the length of the cigar, through which smoke is drawn after the cigar is lit. A cigar rolled with insufficient air passage is referred to by a smoker as "too tight"; one with excessive airflow creating an excessively fast, hot burn is regarded as "too loose". Considerable skill and dexterity on the part of the cigar roller is needed to avoid these opposing pitfalls—a primary factor in the superiority of hand-rolled cigars over their machine-made counterparts. By blending various varieties of filler tobacco, cigar makers create distinctive strength, aroma, and flavor profiles for their various branded products. In general, fatter cigars hold more filler leaves, allowing a greater potential for the creation of complex flavors. In addition to the variety of tobacco employed, the country of origin can be one important determinant of taste, with different growing environments producing distinctive flavors. The fermentation and aging process adds to this variety, as does the particular part of the tobacco plant harvested, with bottom leaves (Spanish: ) having a mild flavor and burning easily, middle leaves (Spanish: ) having a somewhat stronger flavor, with potent and spicy ligero leaves taken from the sun-drenched top of the plant. When used, ligero is always folded into the middle of the filler bunch due to its slow-burning characteristics. Some cigar manufacturers purposely place different types of tobacco from one end to the other to give the cigar smokers a variety of tastes, body, and strength from start to finish. If full leaves are used as filler, a cigar is said to be composed of "long filler". Cigars made from smaller bits of leaf, including many machine-made cigars, are said to be made of "short filler". If a cigar is completely constructed (filler, binder, and wrapper) of tobacco produced in only one country, it is referred to in the cigar industry as a "puro", from the Spanish word for "pure". Size and shape Cigars are commonly categorized by their size and shape, which together are known as the vitola. The size of a cigar is measured by two dimensions: its ring gauge (its diameter in sixty-fourths of an inch) and its length (in inches). In Cuba, next to Havana, there is a display of the world's longest rolled cigars. Parejo The most common shape is the parejo, sometimes referred to as simply "coronas", which have traditionally been the benchmark against which all other cigar formats are measured. They have a cylindrical shape their entire length, one end open, and a round tobacco-leaf "cap" on the other end that must be sliced off, notched, or pierced before smoking. Parejos are designated by the following terms: These dimensions are, at best, idealized. Actual dimensions can vary considerably. Figurado Irregularly shaped cigars are known as figurados and are often priced higher than generally similar sized parejos of a like combination of tobaccos because they are more difficult to make. Historically, especially during the 19th century, figurados were the most popular shapes, but by the 1930s they had fallen out of fashion and all but disappeared. They have recently received a small resurgence in popularity, and currently many manufacturers produce figurados alongside the simpler parejos. The Cuban cigar brand Cuaba only has figurados in their range. Figurados include the following: In practice, the terms Torpedo and Pyramid are often used interchangeably, even among knowledgeable cigar smokers. Min Ron Nee, the Hong Kong-based cigar expert whose work An Illustrated Encyclopaedia of Post-Revolution Havana Cigars is generally considered to be the definitive work on cigars and cigar terms, defines Torpedo as "cigar slang". Nee regards the majority usage of torpedoes as pyramids by another name as acceptable. Arturo Fuente, a large cigar manufacturer based in the Dominican Republic, has also manufactured figurados in exotic shapes ranging from chili peppers to baseball bats and American footballs. They are highly collectible and extremely expensive, when available to the public. Cigarillo A cigarillo is a machine-made cigar that is shorter and narrower than a traditional cigar but larger than little cigars, filtered cigars, and cigarettes, thus similar in size and composition to small panatela sized cigars, cheroots, and traditional blunts. Cigarillos are usually not filtered, although some have plastic or wood tips, and unlike other cigars, some are inhaled when used. Cigarillos are sold in varying quantities: singles, two-packs, three-packs, and five-packs. Cigarillos are very inexpensive: in the United States, usually sold for less than a dollar. Sometimes they are informally called small cigars, mini cigars, or club cigars. Some famous cigar brands, such as Cohiba or Davidoff, also make cigarillos—Cohiba Mini and Davidoff Club Cigarillos, for example. And there are purely cigarillo brands, such as Café Crème, Dannemann Moods, Mehari's, Al Capone, and Swisher Sweets. Cigarillos are often used in making marijuana cigars. Little cigars Little cigars (sometimes called small cigars or miniatures in the UK) differ greatly from regular cigars. They weigh less than cigars and cigarillos, but, more importantly, they resemble cigarettes in size, shape, packaging, and filters. Sales of little cigars quadrupled in the US from 1971 to 1973 in response to the Public Health Cigarette Smoking Act, which banned the broadcast of cigarette advertisements and required stronger health warnings on cigarette packs. Cigars were exempt from the ban, and perhaps more importantly, were taxed at a far lower rate. Little cigars are sometimes called "cigarettes in disguise", and unsuccessful attempts have been made to reclassify them as cigarettes. In the US, sales of little cigars reached an all-time high in 2006, fueled in great part by favorable taxation. In some states, little cigars have successfully been taxed at the rate of cigarettes, such as Illinois, as well as other states. This has caused yet another loophole, in which manufacturers classify their products as "filtered cigars" instead to avoid the higher tax rate. Yet, many continue to argue that there is in fact a distinction between little cigars and filtered cigars. Little cigars offer a similar draw and overall feel to cigarettes, but with aged and fermented tobaccos, while filtered cigars are said to be more closely related to traditional cigars, and are not meant to be inhaled. Research shows that people do inhale smoke from little cigars. Cannagar Recently, with the changing Legality of cannabis, some suppliers are creating so-called "cannagars" (a portmanteau of "cannabis" and "cigar"). These are different from cannabis blunts. Modeled after a traditional cigar, a cannagar is cannabis wrapped within either cannabis or hemp leaf, like a traditional cigar is tobacco wrapped inside dried tobacco leaf. Unlike a cigar, cannagars do not usually contain tobacco, but they do need to be cut and lit like a cigar. Smoking Most machine-made cigars have pre-formed holes in one end or a wood or plastic tip for drawing in the smoke. Hand-rolled cigars require the blunt end to be pierced before lighting. The usual way to smoke a cigar is to not inhale, but to draw the smoke into the mouth. Some smokers inhale the smoke into the lungs, particularly with little cigars. A smoker may swirl the smoke around in the mouth before exhaling it, and may exhale part of the smoke through the nose in order to smell the cigar better as well as to taste it. Cutting Although a handful of cigars are cut or twirled on both ends, the vast majority come with one straight cut end and the other capped with one or more small pieces of wrapper adhered with either a natural tobacco paste or with a mixture of flour and water. The cap end of a cigar must be cut or pierced for the cigar smoke to be drawn properly. The basic types of cigar cutter include: Guillotine (straight cut) Punch V-cut (a.k.a. notch cut, cat's eye, wedge cut, English cut) Grip cutters Cigar Scissors Lighting The head, or cap, of the cigar is usually the end closest to the cigar band, the other the "foot". The band identifies the type of the cigar and may be removed or left on. The smoker cuts or pierces the cap before lighting. The cigar should be rotated during lighting to achieve an even burn while slowly drawn with gentle puffs. If a match is used it should be allowed to burn past its head before being put to the cigar, to avoid imparting unwelcome flavors or chemicals to the smoke. Many specialized gas and fluid lighters are made for lighting cigars. The tip of the cigar should minimally touch any flame, with special care used with torch lighters to avoid charring the tobacco leaves. A third and most traditional way to light a cigar is to use a splinter of cedar known as a spill, which is lit separately before using. Some cigars come individually wrapped in thin cedar sleeves or envelopes, and these can be used to assist in lighting them. Flavor Each brand and type of cigar has its unique taste. Whether a cigar is mild, medium, or full bodied does not correlate with quality. Among the factors which contribute to the scent and flavor of cigar smoke are tobacco types and qualities used for filler, binder, and wrapper, age and aging method, humidity, production techniques (handmade vs. machine-made), and added flavors. Among wrappers, darker tend to produce a sweetness, while lighter usually have a "drier", more neutral taste. Evaluating the flavor of cigars is in some respects similar to wine-tasting. Journals are available for recording personal ratings, description of flavors observed, sizes, brands, etc. Some words used to describe cigar flavor and texture include; spicy, peppery (red or black), sweet, harsh, burnt, green, earthy, woody, cocoa, chestnut, roasted, aged, nutty, creamy, cedar, oak, chewy, fruity, and leathery. Smoke Smoke is produced by incomplete combustion of tobacco during which at least three kinds of chemical reactions occur: pyrolysis breaks down organic molecules into simpler ones, pyrosynthesis recombines these newly formed fragments into chemicals not originally present, and distillation moves compounds such as nicotine from the tobacco into the smoke. For every gram of tobacco smoked, a cigar emits about 120–140 mg of carbon dioxide, 40–60 mg of carbon monoxide, 3–4 mg of isoprene, 1 mg each of hydrogen cyanide and acetaldehyde, and smaller quantities of a large spectrum of volatile N-nitrosamines and volatile organic compounds, with the detailed composition unknown. The most odorous chemicals in cigar smoke are pyridines. Along with pyrazines, they are also the most odorous chemicals in cigar smokers' breath. These substances are noticeable even at extremely low concentrations of a few parts per billion. During smoking, it is not known whether these chemicals are generated by splitting the chemical bonds of nicotine or by Maillard reaction between amino acids and sugars in the tobacco. Cigar smoke is more alkaline than cigarette smoke, and is absorbed more readily by the mucous membrane of the mouth, making it easier for the smoker to absorb nicotine without having to inhale. A single premium cigar may contain as much nicotine as a pack of cigarettes. Parasites Cigars, alongside other tobacco products, can be infested by parasites such as the Lasioderma serricorne (tobacco beetle) and the Ephestia elutella (tobacco moth), which are the most widespread and damaging parasites to the tobacco industry. Infestation can range from the tobacco cultivated in the fields to the leaves used for manufacturing cigars, cigarillos, cigarettes, etc. Both the larvae of Lasioderma serricorne and caterpillars of Ephestia elutella are considered a pest. Humidors The level of humidity in which cigars are kept has a significant effect on their taste and evenness of burn. It is believed that a cigar's flavor best evolves when stored at a relative humidity similar to where the tobacco is grown, and in most cases, the cigars rolled, of approximately 65–70% and a temperature of . Dry cigars become fragile and burn faster while damp cigars burn unevenly and take on a heavy acidic flavor. Humidors are used to maintain an even humidity level. Without one, cigars will lose moisture and acquire the ambient humidity within 2 to 3 days. A humidor's interior lining is typically constructed with three types of wood: Spanish cedar, American (or Canadian) red cedar, and Honduran mahogany. Other materials used for making or lining a humidor are acrylic, tin (mainly seen in older early humidors) and copper, used widely in the 1920s–1950s. Most humidors come with a plastic or metal case with a sponge that works as the humidifier, although most recent versions are of polymer acryl. The latter are filled only with distilled water; the former may use a solution of propylene glycol and distilled water. Humidifiers, and the cigars within them, may become contaminated with bacteria if they are kept too moist. New technologies employing plastic beads or gels which stabilize humidity are becoming widely available. A new humidor requires seasoning, after which a constant humidity must be maintained. The thicker the cedar lining the better. Many humidors contain an analog or digital hygrometer to aid in maintaining a desired humidity level. There are three types of analog: metal spring, natural hair, and synthetic hair. In recent times Electric Humidors, which feature a thermoelectric humidification system have become popular for larger cigar collections. Accessories A wide variety of cigar accessories are available, in varying qualities. Travel case Travel cases protect cigars from direct exposure to the elements and minimize potential damage. Most come in expandable or sturdy leather, although metal leather and plastic lined cases are found. Some feature cardboard or metal tubes for additional protection. Tube Cigar Tubes are used to carry small numbers of cigars, typically one or five, referred to by their number of "fingers". They are usually made from stainless steel, and used for short durations. For longer, a built in humidifier and hygrometer is used. Ashtray Ashtrays are used for collecting the ash produced by the cigar. Such ashtrays are typically larger than those used for cigarette smoking. Holder A cigar holder is small tube in which the end of the cigar is held while smoked, to protect the hand from acquiring the odor of a burning cigar, historically used by women (for cigarettes as well). A cigar stand is a device used to keep a lit cigar out of an ashtray. Health effects Like other forms of tobacco use, cigar smoking poses a significant health risk depending on dosage: risks are greater for those who smoke more cigars, smoke them longer, or inhale more. A review of 22 studies found that cigar smoking is associated with lung cancer, oral cancer, esophageal cancer, pancreatic cancer, oropharyngeal cancer, laryngeal cancer, coronary heart disease (CHD), and aortic aneurysm. Among cigar smokers who reported that they did not inhale, relative mortality (likelihood of death) risk was still highly elevated for oral, esophageal, and laryngeal cancers. Danger of mortality increases proportionally to use, with smokers of one to two cigars per day showing a 2% increase in death rate, compared to non-smokers. The precise statistical health risks to those who smoke less than daily is not established. The depth of inhalation of cigar smoke into the lungs appears to be an important determinant of lung cancer risk: When cigar smokers don't inhale or smoke few cigars per day, the risks are only slightly above those of never smokers. Risks of lung cancer increase with increasing inhalation and with increasing number of cigars smoked per day, but the effect of inhalation is more powerful than that for number of cigars per day. When 5 or more cigars are smoked per day and there is moderate inhalation, the lung cancer risks of cigar smoking approximate those of a one pack per day cigarette smoker. As the tobacco smoke exposure of the lung in cigar smokers increases to approximate the frequency of smoking and depth of inhalation found in cigarette smokers, the difference in lung cancer risks produced by these two behaviors disappears. Cigar smoking can lead to nicotine addiction and cigarette usage. For those who inhale and smoke several cigars a day, the health risk is similar to cigarette smokers. Cigar smoking can also increase the risk of chronic obstructive pulmonary disease (COPD). So-called "little cigars" are commonly inhaled and likely pose the same health risks as cigarettes, while premium cigars are not commonly inhaled or habitually used. Popularity The prevalence of cigar smoking varies depending on location, historical period, and population surveyed. The United States is the top consuming country by total sales by a considerable margin, followed by Germany and the United Kingdom. The U.S. and Western European countries account for about 75% of cigar sales worldwide. United States Consumption of cigars in the U.S. rose from 6.2 billion in 2000 to the peak of an enormous "cigar boom" of 13.8 billion in 2012, which had receded to 11.4 billion by 2015. Among US adults ages 18 and older, 3% reported that they smoke cigars some days or every day (6% of men, 1% of women) in the 2015 National Health Interview Survey. Cigar use among youth declined sharply from 12% reporting having smoked a cigar within the past 30 days approaching the peak of the cigar boom in 2011 to 8% by 2016. Among high school students, cigar use is more common among males (10%) than females (6%). For African American high school students, cigar use is more prevalent (10%) than cigarette use (4%). In popular culture In a reversal of previous decades' portrayal, beginning in the 1980s and 1990s major U.S. print media began to feature cigars favorably. Cigar use was generally framed as a lucrative business or trendy habit, rather than as a major health risk. It is an item whose highest quality is still something most can afford, at least for special occasions. Historic portrayals of the wealthy often caricatured cigar smokers as wearing top hats and tailcoats. Cigars are often given out and smoked to celebrate special occasions, such as the birth of a baby, but also graduations, promotions, and other totems of success. The expression "close but no cigar" comes from the practice of giving away cigars as prizes in fairground games which require the player to hit a target (e.g., a bullseye).	Biology and health sciences	Health and fitness: General	Health
38327	https://en.wikipedia.org/wiki/Cigarette	Cigarette	A cigarette is a narrow cylinder containing a combustible material, typically tobacco, that is rolled into thin paper for smoking. The cigarette is ignited at one end, causing it to smolder; the resulting smoke is orally inhaled via the opposite end. Cigarette smoking is the most common method of tobacco consumption. The term cigarette, as commonly used, refers to a tobacco cigarette, but the word is sometimes used to refer to other substances, such as a cannabis cigarette or a herbal cigarette. A cigarette is distinguished from a cigar by its usually smaller size, use of processed leaf, different smoking method, and paper wrapping, which is typically white. There are significant negative health effects from smoking cigarettes such as cancer, chronic obstructive pulmonary disease (COPD), heart disease, birth defects, and other health problems relating to nearly every organ of the body. Most modern cigarettes are filtered, although this does not make the smoke inhaled from them contain fewer carcinogens and harmful chemicals. Nicotine, the psychoactive drug in tobacco, makes cigarettes highly addictive. About half of cigarette smokers die of tobacco-related disease and lose on average 14 years of life. Every year, cigarette smoking causes more than 8 million deaths worldwide; more than 1.3 million of these are non-smokers dying as the result of exposure to secondhand smoke. These harmful effects have led to legislation that has prohibited smoking in many workplaces and public areas, regulated marketing and purchasing age of tobacco, and levied taxes to discourage cigarette use. In the 21st century electronic cigarettes (also called e-cigarettes or vapes) were developed, whereby a substance contained within (typically a liquid solution containing nicotine) is vaporized by a battery-powered heating element as opposed to being burned. Such devices are commonly promoted by their manufacturers as safer alternatives to conventional cigarettes. Since e-cigarettes are a relatively new product, scientists do not have data on their possible long-term health effects, but there are significant health risks associated with their use. History Global The earliest forms of cigarettes were similar to their predecessor, the cigar. Cigarettes appear to have had antecedents in Mexico and Central America around the 9th century in the form of reeds and smoking tubes. The Maya, and later the Aztecs, smoked tobacco and other psychoactive drugs in religious rituals and frequently depicted priests and deities smoking on pottery and temple engravings. The cigarette and the cigar were the most common methods of smoking in the Caribbean, Mexico, and Central and South America until recent times. The North American, Central American, and South American cigarette used various plant wrappers; when it was brought back to Spain, maize wrappers were introduced, and by the 17th century, fine paper. The resulting product was called papelate and is documented in Goya's paintings La Cometa, La Merienda en el Manzanares, and El juego de la pelota a pala (18th century). By 1830 the cigarette had become known in France, where it received the name cigarette, and in 1845 the French state tobacco monopoly began manufacturing them. The French word made its way into English in the 1840s. Some American reformers promoted the spelling cigaret, but this was never widespread and is now largely abandoned. The first patented cigarette-making machine was invented by Juan Nepomuceno Adorno of Mexico in 1847. In the 1850s, Turkish cigarette leaves had become popular. However, production climbed markedly when another cigarette-making machine was developed in the 1880s by James Albert Bonsack, which vastly increased the productivity of cigarette companies, which went from making about 40,000 hand-rolled cigarettes daily to around 4 million. At the time, these imported cigarettes from the United States had significant sales among British smokers. In the English-speaking world, the use of tobacco in cigarette form became increasingly widespread during and after the Crimean War, when British soldiers began emulating their Ottoman Turkish comrades and Russian enemies, who had begun rolling and smoking tobacco in strips of old newspaper for lack of proper cigar-rolling leaf. This was helped by the development of tobaccos suitable for cigarette use, and by the development of the Egyptian cigarette export industry. Cigarettes may have been initially used in a manner similar to pipes, cigars, and cigarillos and not inhaled. As cigarette tobacco became milder and more acidic, inhaling may have become perceived as more agreeable; a sentiment supported by advertising in the 1930s. However, Helmuth von Moltke noticed in the 1830s that Ottomans (and he himself) inhaled the Turkish tobacco and Latakia from their pipes (which are both initially sun-cured, acidic leaf varieties). The widespread smoking of cigarettes in the Western world is largely a 20th-century phenomenon. By the late 19th century cigarettes were known as coffin nails but the link between lung cancer and smoking was not established until the 20th century. German doctors were the first to make the link, and it led to the first antitobacco movement in Nazi Germany. During World War I and World War II, cigarettes were rationed to soldiers. During the Vietnam War, cigarettes were included with C-ration meals. In 1975, the U.S. government stopped putting cigarettes in military rations. During the second half of the 20th century, the adverse health effects of tobacco smoking started to become widely known and printed health warnings became common on cigarette packets. Graphical cigarette warning labels are a more effective method to communicate to the public the dangers of cigarette smoking. Canada, Mexico, Belgium, Denmark, Sweden, Thailand, Malaysia, India, Pakistan, Australia, Argentina, Brazil, Chile, Peru, Greece, the Netherlands, New Zealand, Norway, Hungary, the United Kingdom, France, Romania, Singapore, Egypt, Jordan, Nepal and Turkey all have both textual warnings and graphic visual images displaying, among other things, the damaging effects tobacco use has on the human body. The United States has implemented textual but not graphical warnings. The cigarette has evolved much since its conception; for example, the thin bands that travel transverse to the "axis of smoking" (thus forming circles along the length of the cigarette) are alternate sections of thin and thick paper to facilitate effective burning when being drawn, and retard burning when at rest. Synthetic particulate filters may remove some of the tar before it reaches the smoker. The "holy grail" for cigarette companies has been a cancer-free cigarette. On record, the closest historical attempt was produced by scientist James Mold. Under the name project TAME, he produced the XA cigarette. However, in 1978, his project was terminated. Since 1950, the average nicotine and tar content of cigarettes has steadily fallen. Research has shown that the fall in overall nicotine content has led to smokers inhaling larger volumes per puff. United States One entrepreneur who was quick to spot the advantages of machine-made cigarettes was James Buchanan Duke. Previously a producer of smoking tobacco only, his firm, W. Duke & Sons & Co., entered the cigarette industry in the early 1880s. After installing two Bonsack machines, Duke spent heavily on advertising and sales promotion with the result that by 1889 his was the largest cigarette manufacturer in the country. The new Bonsack machines were of decisive importance in rapid, cheap manufacture of all tobacco products but one. Cigars needed slow, laborious hand rolling and were produced in hundreds of small workshops, especially in New York City. In 1890 Duke and the other four major cigarette companies combined to form the American Tobacco Company, a firm that dominated the market and used aggressive tactics on hundreds of small competitors until they sold out. It was called the "Tobacco Trust." The trust soon expanded its operations to include cigars, smoking, chewing tobacco and snuff. Among the companies drawn into this organization were plug manufacturers, Liggett & Myers and R. J. Reynolds Tobacco Company, which at the time produced twist and flat plug, and P. Lorillard, an old-line manufacturer of snuff. By 1910 the trust produced 86% of all cigarettes produced in the United States, and 75% to 95% of other forms, but only 14% of the cigars. At the start of the 20th century, the per capita annual consumption in the U.S. was 54 cigarettes (with fewer than 0.5% of the population smoking more than 100 cigarettes per year), and consumption there peaked at 4,259 per capita in 1965. At that time, about 50% of men and 33% of women smoked (defined as smoking more than 100 cigarettes per year). By 2000, consumption had fallen to 2,092 per capita, corresponding to about 30% of men and 22% of women smoking more than 100 cigarettes per year, and by 2006 per capita consumption had declined to 1,691; implying that about 21% of the population smoked 100 cigarettes or more per year. Construction Manufacturers have described the cigarette as "a drug administration system for the delivery of nicotine in acceptable and attractive form". Modern commercially manufactured cigarettes are seemingly simple objects consisting mainly of a tobacco blend, paper, PVA glue to bond the outer layer of paper together, and often also a cellulose acetate–based filter. While the assembly of cigarettes is straightforward, much focus is given to the creation of each of the components, in particular the tobacco blend. A key ingredient that makes cigarettes more addictive is the inclusion of reconstituted tobacco, which has additives to make nicotine more volatile as the cigarette burns. Paper The paper for holding the tobacco blend may vary in porosity to allow ventilation of the burning ember or contain materials that control the burning rate of the cigarette and stability of the produced ash. The papers used in tipping the cigarette (forming the mouthpiece) and surrounding the filter stabilize the mouthpiece from saliva and moderate the burning of the cigarette, as well as the delivery of smoke with the presence of one or two rows of small laser-drilled air holes. Tobacco blend The process of blending gives the end product a consistent taste from batches of tobacco grown in different areas of a country that may change in flavor profile from year to year due to different environmental conditions. Modern cigarettes produced after the 1950s, although composed mainly of shredded tobacco leaf, use a significant quantity of tobacco processing byproducts in the blend. Each cigarette's tobacco blend is made mainly from the leaves of flue-cured brightleaf, burley tobacco, and oriental tobacco. These leaves are selected, processed, and aged prior to blending and filling. The processing of brightleaf and burley tobaccos for tobacco leaf "strips" produces several byproducts such as leaf stems, tobacco dust, and tobacco leaf pieces ("small laminate"). To improve the economics of producing cigarettes, these byproducts are processed separately into forms where they can then be added back into the cigarette blend without an apparent or marked change in the cigarette's quality. The most common tobacco byproducts include: Blended leaf (BL) sheet: a thin, dry sheet cast from a paste made with tobacco dust collected from tobacco stemming, finely milled burley-leaf stem, and pectin. Reconstituted leaf (RL) sheet: a paper-like material made from recycled tobacco fines, tobacco stems and "class tobacco", which consists of tobacco particles less than 30 mesh in size (about 0.6 mm) that are collected at any stage of tobacco processing: RL is made by extracting the soluble chemicals in the tobacco byproducts, processing the leftover tobacco fibers from the extraction into a paper, and then reapplying the extracted materials in concentrated form onto the paper in a fashion similar to what is done in paper sizing. At this stage, ammonium additives are applied to make reconstituted tobacco an effective nicotine delivery system. Expanded (ES) or improved stem (IS): ES is rolled, flattened, and shredded leaf stems that are expanded by being soaked in water and rapidly heated. Improved stem follows the same process, but is simply steamed after shredding. Both products are then dried. These products look similar in appearance, but are different in taste. According to data from the World Health Organization, the amount of tobacco per 1000 cigarettes fell from in 1960 to in 1999, largely as a result of reconstituting tobacco, fluffing, and additives. A recipe-specified combination of brightleaf, burley-leaf, and oriental-leaf tobacco is mixed with various additives to improve its flavors. Most commercially available cigarettes today contain tobacco that is treated with sugar to counter the harshness of the smoke. Additives Various additives are combined into the shredded tobacco product mixtures, with humectants such as propylene glycol or glycerol, as well as flavoring products and enhancers such as cocoa solids, licorice, tobacco extracts, and various sugars, which are known collectively as "casings". The leaf tobacco is then shredded, along with a specified amount of small laminate, expanded tobacco, BL, RL, ES, and IS. A perfume-like flavor/fragrance, called the "topping" or "toppings", which is most often formulated by flavor companies, is then blended into the tobacco mixture to improve the consistency in flavor and taste of the cigarettes associated with a certain brand name. Additionally, they replace lost flavors due to the repeated wetting and drying used in processing the tobacco. Finally, the tobacco mixture is filled into cigarette tubes and packaged. A list of 599 cigarette additives, created by five major American cigarette companies, was approved by the Department of Health and Human Services in April 1994. None of these additives is listed as an ingredient on the cigarette packs. Chemicals are added for organoleptic purposes and many boost the addictive properties of cigarettes, especially when burned. One of the classes of chemicals on the list, ammonia salts, convert bound nicotine molecules in tobacco smoke into free nicotine molecules. This process, known as freebasing, could potentially increase the effect of nicotine on the smoker, but experimental data suggests that absorption is, in practice, unaffected. Cigarette tube Cigarette tubes are prerolled cigarette paper usually with an acetate or paper filter at the end. They have an appearance similar to a finished cigarette, but are without any tobacco or smoking material inside. The length varies from Regular (70 mm) to King Size (84 mm) as well as 100s (100 mm) and 120s (120 mm). Filling a cigarette tube is usually done with a cigarette injector (also known as a shooter). Cone-shaped cigarette tubes, known as cones, can be filled using a packing stick or straw because of their shape. Cone smoking is popular because as the cigarette burns, it tends to get stronger and stronger. A cone allows more tobacco to be burned at the beginning than the end, allowing for an even flavor The United States Tobacco Taxation Bureau defines a cigarette tube as "Cigarette paper made into a hollow cylinder for use in making cigarettes." Cigarette filter A cigarette filter or filter tip is a component of a cigarette. Filters are typically made from cellulose acetate fibre. Most factory-made cigarettes are equipped with a filter; those who roll their own can buy them separately. Filters can reduce some substances from smoke but do not make cigarettes any safer to smoke. Cigarette butt In North America, the common name for the remains of a cigarette after smoking is a cigarette butt. In Britain, it is also called a fag-end or a dog-end. The butt is typically about 30% of the cigarette's original length. It consists of a tissue tube which holds a filter and some remains of tobacco mixed with ash. They are the most numerically frequent litter in the world. Cigarette butts accumulate outside buildings, on parking lots, and streets where they can be transported through storm drains to streams, rivers, and beaches. In a 2013 trial, the city of Vancouver, British Columbia, partnered with TerraCycle to create a system for recycling of cigarette butts. A reward of 1¢ per collected butt was offered to determine the effectiveness of a deposit system similar to that of beverage containers. Electronic cigarette An electronic cigarette (commonly known as a vape) is a handheld battery-powered vaporizer that simulates smoking by providing some of the behavioral aspects of smoking, including the hand-to-mouth action of smoking, but without combusting tobacco. Using an e-cigarette is known as "vaping" and the user is referred to as a "vaper". Instead of cigarette smoke, the user inhales an aerosol, commonly called vapor. E-cigarettes typically have a heating element that atomizes a liquid solution called e-liquid. E-cigarettes are automatically activated by taking a puff; others turn on manually by pressing a button. Some e-cigarettes look like traditional cigarettes, but they come in many variations. Most versions are reusable, though some are disposable. There are first-generation, second-generation, third-generation, and fourth-generation devices. E-liquids usually contain propylene glycol, glycerin, nicotine, flavorings, additives, and differing amounts of contaminants. E-liquids are also sold without propylene glycol, nicotine, or flavors. The benefits and the health risks of e-cigarettes are uncertain. There is moderate-certainty evidence that e-cigarettes with nicotine may help people quit smoking when compared with e-cigarettes without nicotine and nicotine replacement therapy. However, other studies have not supported the finding that e-cigarettes are proven to be more effective than smoking cessation medicine. There is concern with the possibility that non-smokers and children may start nicotine use with e-cigarettes at a rate higher than anticipated than if they were never created. Following the possibility of nicotine addiction from e-cigarette use, there is concern children may start smoking cigarettes. Youth who use e-cigarettes are more likely to go on to smoke cigarettes. Their part in tobacco harm reduction is unclear, while another review found they appear to have the potential to lower tobacco-related death and disease. Regulated US Food and Drug Administration nicotine replacement products may be safer than e-cigarettes, but e-cigarettes are generally seen as safer than combusted tobacco products. It is estimated their safety risk to users is similar to that of smokeless tobacco. The long-term effects of e-cigarette use are unknown. The risk from serious adverse events was reported in 2016 to be low. Less serious adverse effects include abdominal pain, headache, blurry vision, throat and mouth irritation, vomiting, nausea, and coughing. Nicotine itself is associated with some health harms. In 2019 and 2020, an outbreak of severe lung illness throughout the US was linked to the use of vaping products E-cigarettes create vapor made of fine and ultrafine particles of particulate matter, which have been found to contain propylene glycol, glycerin, nicotine, flavors, small amounts of toxicants, carcinogens, and heavy metals, as well as metal nanoparticles, and other substances. Its exact composition varies across and within manufacturers, and depends on the contents of the liquid, the physical and electrical design of the device, and user behavior, among other factors. E-cigarette vapor potentially contains harmful chemicals not found in tobacco smoke. E-cigarette vapor contains fewer toxic chemicals, and lower concentrations of potential toxic chemicals than cigarette smoke. The vapor is probably much less harmful to users and bystanders than cigarette smoke, although concern exists that the exhaled vapor may be inhaled by non-users, particularly indoors. Health effects Smokers The harm from smoking comes from the many toxic chemicals in the natural tobacco leaf and those formed in smoke from burning tobacco. People keep smoking because the nicotine, the primary psychoactive chemical in cigarettes, is highly addictive. Cigarettes, like narcotics, have been described as "strategically addictive", with the addictive properties being a core component of the business strategy. About half of smokers die from a smoking-related cause. Smoking harms nearly every organ of the body. Smoking leads most commonly to diseases affecting the heart, liver, and lungs, being a major risk factor for heart attacks, strokes, chronic obstructive pulmonary disease (COPD) (including emphysema and chronic bronchitis), and cancer (particularly lung cancer, cancers of the larynx and mouth, and pancreatic cancer). It also causes peripheral vascular disease and hypertension. Children born to women who smoke during pregnancy are at higher risk of congenital disorders, cancer, respiratory disease, and sudden death. On average, each cigarette smoked is estimated to shorten life by 11 minutes. Starting smoking earlier in life and smoking cigarettes higher in tar increases the risk of these diseases. The World Health Organization estimates that tobacco causes 8 million deaths each year as of 2019 and 100 million deaths over the course of the 20th century. Cigarettes produce an aerosol containing over 4,000 chemical compounds, including nicotine, carbon monoxide, acrolein, and oxidant substances. Over 70 of these are carcinogens. The most important chemical compounds causing cancer are those that produce DNA damage since such damage appears to be the primary underlying cause of cancer. Cunningham et al. combined the microgram weight of the compound in the smoke of one cigarette with the known genotoxic effect per microgram to identify the most carcinogenic compounds in cigarette smoke. The seven most important carcinogens in tobacco smoke are shown in the table, along with DNA alterations they cause. "Ulcerative colitis is a condition of nonsmokers in which nicotine is of therapeutic benefit." A recent review of the available scientific literature concluded that the apparent decrease in Alzheimer disease risk may be simply because smokers tend to die before reaching the age at which it normally occurs. "Differential mortality is always likely to be a problem where there is a need to investigate the effects of smoking in a disorder with very low incidence rates before age 75 years, which is the case of Alzheimer's disease", it stated, noting that smokers are only half as likely as nonsmokers to survive to the age of 80. Gateway theory A very strong argument has been made about the association between adolescent exposure to nicotine by smoking conventional cigarettes and the subsequent onset of using other dependence-producing substances. Strong temporal and dose-dependent associations have been reported, and a plausible biological mechanism (via rodent and human modeling) suggests that long-term changes in the neural reward system take place as a result of adolescent smoking. Adolescent smokers of conventional cigarettes have disproportionately high rates of comorbid substance use, and longitudinal studies have suggested that early adolescent smoking may be a starting point or "gateway" for substance use later in life, with this effect more likely for persons with attention deficit hyperactivity disorder (ADHD). Although factors such as genetic comorbidity, innate propensity for risk-taking, and social influences may underlie these findings, both human neuroimaging and animal studies suggest a neurobiological mechanism also plays a role. In addition, behavioral studies in adolescent and young adult smokers have revealed an increased propensity for risk-taking, both generally and in the presence of peers, and neuroimaging studies have shown altered frontal neural activation during a risk-taking task as compared with nonsmokers. In 2011, Rubinstein and colleagues used neuroimaging to show decreased brain response to a natural reinforcer (pleasurable food cues) in adolescent light smokers (1–5 cigarettes per day), with their results highlighting the possibility of neural alterations consistent with nicotine dependence and altered brain response to reward even in adolescent low-level smokers. Second-hand smoke Second-hand smoke is a mixture of smoke from the burning end of a cigarette and the smoke exhaled from the lungs of smokers. It is involuntarily inhaled, lingers in the air for hours after cigarettes have been extinguished, and can cause a wide range of adverse health effects, including cancer, respiratory infections, and asthma. Nonsmokers who are exposed to second-hand smoke at home or work increase their heart disease risk by 25–30% and their lung cancer risk by 20–30%. Second-hand smoke has been estimated to cause 38,000 deaths per year, of which 3,400 are deaths from lung cancer in nonsmokers. Sudden infant death syndrome, ear infections, respiratory infections, and asthma attacks can occur in children who are exposed to second-hand smoke. Scientific evidence shows that no level of exposure to second-hand smoke is safe. Legislation Smoking restrictions Many governments impose restrictions on smoking tobacco, especially in public areas. The primary justification has been the negative health effects of second-hand smoke. Laws vary by country and locality. Nearly all countries have laws restricting places where people can smoke in public, and over 40 countries have comprehensive smoke-free laws that prohibit smoking in virtually all public venues. Smoking age In the United States, the age to buy tobacco products is 21 in all states as of 2020. Similar laws exist in many other countries. In Canada, most of the provinces require smokers to be 19 years of age to purchase cigarettes (except for Quebec and the prairie provinces, where the age is 18). However, the minimum age only concerns the purchase of tobacco, not use. Alberta, however, does have a law which prohibits the possession or use of tobacco products by all persons under 18, punishable by a $100 fine. Australia, New Zealand, Poland, and Pakistan have a nationwide ban on the selling of all tobacco products to people under the age of 18. Since October 1, 2007, it has been illegal for retailers to sell tobacco in all forms to people under the age of 18 in three of the UK's four constituent countries (England, Wales, Northern Ireland, and Scotland), rising from 16. It is also illegal to sell lighters, rolling papers, and all other tobacco-associated items to people under 18. It is not illegal for people under 18 to buy or smoke tobacco, just as it was not previously for people under 16; it is only illegal for the said retailer to sell the item. The age increase from 16 to 18 came into force in Northern Ireland on September 1, 2008. In the Republic of Ireland, bans on the sale of the smaller 10-packs and confectionery that resembles tobacco products (candy cigarettes) came into force on May 31, 2007, in a bid to cut underaged smoking. In October 2023 it was announced that the government proposed introducing a ban on sales of cigarettes to anyone born after 2008. Most countries in the world have a legal vending age of 18. In North Macedonia, Italy, Malta, Austria, Luxembourg, and Belgium, the age for legal vending is 16. Since January 1, 2007, all cigarette machines in public places in Germany must attempt to verify a customer's age by requiring the insertion of a debit card. Turkey, which has one of the highest percentage of smokers in its population, has a legal age of 18. Japan is one of the highest tobacco-consuming nations, and requires purchasers to be 20 years of age. Since July 2008, Japan has enforced this age limit at cigarette vending machines through use of the taspo smart card. In other countries, such as Egypt, it is legal to use and purchase tobacco products regardless of age. Germany raised the purchase age from 16 to 18 on September 1, 2007. Some police departments in the United States occasionally send an underaged teenager into a store where cigarettes are sold, and have the teen attempt to purchase cigarettes, with their own or no ID. If the vendor then completes the sale, the store is issued a fine. Similar enforcement practices are regularly performed by Trading Standards officers in the UK, Israel, and the Republic of Ireland. Taxation Cigarettes are taxed both to reduce use, especially among youth, and to raise revenue. Higher prices for cigarettes discourage smoking. Every 10% increase in the price of cigarettes reduces youth smoking by about 7% and overall cigarette consumption by about 4%. The World Health Organization (WHO) recommends that globally cigarettes be taxed at a rate of three-quarters of cigarettes sale price as a way of deterring cancer and other negative health outcomes. Cigarette sales are a significant source of tax revenue in many localities. This fact has historically been an impediment for health groups seeking to discourage cigarette smoking, since governments seek to maximize tax revenues. Furthermore, some countries have made cigarettes a state monopoly, which has the same effect on the attitude of government officials outside the health field. In the United States, states are a primary determinant of the total tax rate on cigarettes. Generally, states that rely on tobacco as a significant farm product tend to tax cigarettes at a low rate. Coupled with the federal cigarette tax of $1.01 per pack, total cigarette-specific taxes range from $1.18 per pack in Missouri to $8.00 per pack in Silver Bay, New York. As part of the Family Smoking Prevention and Tobacco Control Act, the federal government collects user fees to fund Food and Drug Administration (FDA) regulatory measures over tobacco. Fire-safe cigarette Cigarettes are a frequent source of deadly fires in private homes, which prompted both the European Union and the United States to require cigarettes to be fire-standard compliant. According to Simon Chapman, a professor of public health at the University of Sydney, reduction of burning agents in cigarettes would be a simple and effective means of dramatically reducing the ignition propensity of cigarettes. Since the 1980s, prominent cigarette manufacturers such as Philip Morris and R.J. Reynolds have developed fire safe cigarettes, but Phillip Morris was later the subject of a government lawsuit for allegedly hiding the even greater dangers associated with their brand of such cigarettes. The burn rate of cigarette paper is regulated through the application of different forms of microcrystalline cellulose to the paper. Cigarette paper has been specially engineered by creating bands of different porosity to create "fire-safe" cigarettes. These cigarettes have a reduced idle burning speed which allows them to self-extinguish. This fire-safe paper is manufactured by mechanically altering the setting of the paper slurry. New York was the first U.S. state to mandate that all cigarettes manufactured or sold within the state comply with a fire-safe standard. Canada has passed a similar nationwide mandate based on the same standard. All U.S. states are gradually passing fire-safe mandates. The European Union in 2011 banned cigarettes that do not meet a fire-safety standard. According to a study made by the European Union in 16 European countries, 11,000 fires were due to people carelessly handling cigarettes between 2005 and 2007. This caused 520 deaths with 1,600 people injured. Cigarette advertising Many countries have restrictions on cigarette advertising, promotion, sponsorship, and marketing. For example, in the Canadian provinces of British Columbia, Saskatchewan and Alberta, the retail store display of cigarettes is completely prohibited if persons under the legal age of consumption have access to the premises. In Ontario, Manitoba, Newfoundland and Labrador, and Quebec, Canada and the Australian Capital Territory the display of tobacco is prohibited for everyone, regardless of age, as of 2010. This retail display ban includes noncigarette products such as cigars and blunt wraps. Warning messages in packages As a result of tight advertising and marketing prohibitions, tobacco companies look at the pack differently: they view it as a strong component in displaying brand imagery and a creating significant in-store presence at the point of purchase. Market testing shows the influence of this dimension in shifting the consumer's choice when the same product displays in an alternative package. Companies have manipulated a variety of elements in packs designs to communicate the impression of lower in tar or milder cigarettes, whereas the components were the same. Some countries require cigarette packs to contain warnings about health hazards. The United States was the first, later followed by other countries including Canada, most of Europe, Australia, Pakistan, India, Hong Kong, and Singapore. In 1985, Iceland became the first country to enforce graphic warnings on cigarette packaging. At the end of December 2010, new regulations from Ottawa increased the size of tobacco warnings to cover three-quarters of the cigarette package in Canada. As of November 2010, 39 countries have adopted similar legislation. In February 2011, the Canadian government passed regulations requiring cigarette packs to contain 12 new images to cover 75% of the outside panel and eight new health messages on the inside panel with full color. As of April 2011, Australian regulations require all packs to use a bland olive green that researchers determined to be the least attractive color, with 75% coverage on the front of the pack and all of the back consisting of graphic health warnings. The only feature that differentiates one brand from another is the product name in a standard color, position, font size, and style. Similar policies have since been adopted in France and the United Kingdom. In response to these regulations, Philip Morris International, Japan Tobacco Inc., British American Tobacco Plc., and Imperial Tobacco attempted to sue the Australian government. On August 15, 2012, the High Court of Australia dismissed the suit and made Australia the first country to introduce brand-free plain cigarette packaging with health warnings covering 90 and 70% of back and front packaging, respectively. This took effect on December 1, 2012. Prohibition of tobacco A few countries have outlawed tobacco completely or made plans to do so. In 2004, Bhutan became the first country in the world to completely outlaw the cultivation, harvesting, production, and sale of tobacco and tobacco products. Enforcement of the prohibition increased with the passage of the Tobacco Control Act of Bhutan 2010. However, small allowances for personal possession are permitted as long as the possessors can prove that they have paid import duties. The Pitcairn Islands had previously banned the sale of cigarettes, but it now permits sales from a government-run store. The Pacific island of Niue hopes to become the next country to prohibit the sale of tobacco. Iceland is also proposing banning tobacco sales from shops, making it prescription-only and therefore dispensable only in pharmacies on doctor's orders. Singapore and the Australian state of Tasmania have proposed a 'tobacco free millennium generation initiative' by banning the sale of all tobacco products to anyone born in and after the year 2000. In March 2012, Brazil became the world's first country to ban all flavored tobacco including menthols. It also banned the majority of the estimated 600 additives used, permitting only eight. This regulation applies to domestic and imported cigarettes. Tobacco manufacturers had 18 months to remove the noncompliant cigarettes, 24 months to remove the other forms of noncompliant tobacco. Under sharia law, the consumption of cigarettes by Muslims is prohibited. Environmental effects Cigarette filters are made up of thousands of polymer chains of cellulose acetate, which has the chemical structure shown to the right. Once discarded into the environment, the filters create a large waste problem. Cigarette filters are the most common form of litter in the world, as approximately 5.6 trillion cigarettes are smoked every year worldwide. Of those, an estimated 4.5 trillion cigarette filters become litter every year. To develop an idea of the waste weight amount produced a year the table below was created. Discarded cigarette filters usually end up in the water system through drainage ditches and are transported by rivers and other waterways to the ocean. Aquatic life health concerns In the 2006 International Coastal Cleanup, cigarettes and cigarette butts constituted 24.7% of the total collected pieces of garbage, over twice as many as any other category, which is not surprising seeing the numbers in the table above of waste produced each year. Cigarette filters contain the chemicals filtered from cigarettes and can leach into waterways and water supplies. The toxicity of used cigarette filters depends on the specific tobacco blend and additives used by the cigarette companies. After a cigarette is smoked, the filter retains some of the chemicals, and some of those are considered carcinogenic. When studying the environmental effects of cigarette filters, the various chemicals that can be found in cigarette filters are not studied individually, due to the complexity of doing so. Researchers instead focus on the whole cigarette filter and its LD50. LD50 is defined as the lethal dose that kills 50% of a sample population. This allows for a simpler study of the toxicity of cigarette filters. One recent study has looked at the toxicity of smoked cigarette filters (smoked filter + tobacco), smoked cigarette filters (no tobacco), and unsmoked cigarette filters (no tobacco). The results of the study showed that for the LD50 of both marine topsmelt (Atherinops affinis) and freshwater fathead minnow (Pimephales promelas), smoked cigarette filters + tobacco are more toxic than smoked cigarette filters, but both are severely more toxic than unsmoked cigarette filters. Other health concerns Toxic chemicals are not the only human health concern to take into considerations; the others are cellulose acetate and carbon particles that are breathed in while smoking. These particles are suspected of causing lung damage. The next health concern is that of plants. Under certain growing conditions, plants on average grow taller and have longer roots than those exposed to cigarette filters in the soil. A connection exists between cigarette filters introduced to soil and the depletion of some soil nutrients over time. Another health concern to the environment is not only the toxic carcinogens that are harmful to the wildlife, but also the filters themselves pose an ingestion risk to wildlife that may presume filter litter as food. The last major health concern to make note of for marine life is the toxicity that deep marine topsmelt and fathead minnow pose to their predators. This could lead to toxin build-up (bioaccumulation) in the food chain and have long reaching negative effects. Smoldering cigarette filters have also been blamed for triggering fires from residential areas to major wildfires and bushfires which has caused major property damage and also death as well as disruption to services by triggering alarms and warning systems. Degradation Once in the environment, cellulose acetate can go through biodegradation and photodegradation. Several factors go into determining the rate of each degradation process. This variance in rate and resistance to biodegradation in many conditions is a factor in littering and environmental damage. Biodegradation The first step in the biodegradation of cellulose acetate is the deactylation of the acetate from the polymer chain (which is the opposite of acetylation). An acetate is a negative ion with the chemical formula of C2H3O2−. Deacetylation can be performed by either chemical hydrolysis or acetylesterase. Chemical hydrolysis is the cleavage of a chemical bond by addition of water. In the reaction, water (H2O) reacts with the acetic ester functional group attached the cellulose polymer chain and forms an alcohol and acetate. The alcohol is simply the cellulose polymer chain with the acetate replaced with an alcohol group. The second reaction is exactly the same as chemical hydrolysis with the exception of the use of an acetylesterase enzyme. The enzyme, found in most plants, catalyzes the chemical reaction shown below. acetic ester + H2O alcohol + acetate In the case of the enzymatic reaction, the two substrates (reactants) are again acetic ester and H2O, the two products of the reaction are alcohol and acetate. This reaction is exactly the same as the chemical hydrolysis. Both of these products are perfectly fine in the environment. Once the acetate group is removed from the cellulose chain, the polymer can be readily degraded by cellulase, which is another enzyme found in fungi, bacteria, and protozoans. Cellulases break down the cellulose molecule into monosaccharides ("simple sugars") such as beta-glucose, or shorter polysaccharides and oligosaccharides. These simple sugars are not harmful to the environment and are in fact are a useful product for many plants and animals. The breakdown of cellulose is of interest in the field of biofuel. Due to the conditions that affect the process, large variation in the degradation time of cellulose acetate occurs. Factors in biodegradation The duration of the biodegradation process is cited as taking as little as one month to as long as 15 years or more, depending on the environmental conditions. The major factor that affects the biodegradation duration is the availability of acetylesterase and cellulase enzymes. Without these enzymes, biodegradation only occurs through chemical hydrolysis and stops there. Temperature is another major factor: if the organisms that contain the enzymes are too cold to grow, then biodegradation is severely hindered. Availability of oxygen in the environment also affects the degradation. Cellulose acetate is degraded within 2–3 weeks under aerobic assay systems of in vitro enrichment cultivation techniques and an activated sludge wastewater treatment system. It is degraded within 14 weeks under anaerobic conditions of incubation with special cultures of fungi. Ideal conditions were used for the degradation (i.e., a suitable temperature, and available organisms to provide the enzymes). Thus, filters last longer in places with low oxygen concentration, such as swamps and bogs. Overall, the biodegradation process of cellulose acetate is not an instantaneous process. Photodegradation The other process of degradation is photodegradation, which is when a molecular bond is broken by the absorption of photon radiation (i.e. light). Due to cellulose acetate carbonyl groups, the molecule naturally absorbs light at 260 nm, but it contains some impurities which can absorb light. These impurities are known to absorb light in the far UV light region (< 280 nm). The atmosphere filters radiation from the sun and allows radiation of > 300 nm only to reach the surface. Thus, the primary photodegradation of cellulose acetate is considered insignificant to the total degradation process, since cellulose acetate and its impurities absorb light at shorter wavelengths. Research is focused on the secondary mechanisms of photodegradation of cellulose acetate to help make up for some of the limitations of biodegradation. The secondary mechanisms would be the addition of a compound to the filters that would be able to absorb natural light and use it to start the degradation process. The main two areas of research are in photocatalytic oxidation and photosensitized degradation. Photocatalytic oxidation uses a species that absorbs radiation and creates hydroxyl radicals that react with the filters and start the breakdown. Photosensitized degradation, though, uses a species that absorbs radiation and transfers the energy to the cellulose acetate to start the degradation process. Both processes use other species that absorbed light > 300 nm to start the degradation of cellulose acetate. Solution and remediation projects Several options are available to help reduce the environmental effects of cigarette butts. Proper disposal into receptacles leads to decreased numbers found in the environment and their effect on the environment. Another method is making fines and penalties for littering filters; many governments have sanctioned stiff penalties for littering of cigarette filters; for example, Washington imposes a penalty of $1,025 for littering cigarette filters. Another option is developing better biodegradable filters; much of this work relies heavily on the research in the secondary mechanism for photodegradation as stated above, but a new research group has developed an acid tablet that goes inside the filters, and once wet enough, releases acid that speeds up the degradation to around two weeks. The research is still only in test phase and the hope is soon it will go into production. The next option is using cigarette packs with a compartment in which to discard cigarette butts, implementing monetary deposits on filters, increasing the availability of butt receptacles, and expanding public education. It may even be possible to ban the sale of filtered cigarettes altogether on the basis of their adverse environmental effects. Recent research has been put into finding ways to use the filter waste to develop a desired product. One research group in South Korea has developed a simple one-step process that converts the cellulose acetate in discarded cigarette filters into a high-performing material that could be integrated into computers, handheld devices, electrical vehicles, and wind turbines to store energy. These materials have demonstrated superior performance as compared to commercially available carbon, grapheme, and carbon nanotubes. The product is showing high promise as a green alternative for the waste problem. Consumption Smoking has become less popular, but is still a large public health problem globally. Worldwide, smoking rates fell from 41% in 1980 to 31% in 2012, although the actual number of smokers increased because of population growth. In 2017, 5.4 trillion cigarettes were produced globally, and were smoked by almost 1 billion people. Smoking rates have leveled off or declined in most countries, but is increasing in some low- and middle-income countries. The significant reductions in smoking rates in the United States, United Kingdom, Australia, Brazil, and other countries that implemented strong tobacco control programs have been offset by the increasing consumption in low income countries, especially China. The Chinese market now consumes more cigarettes than all other low- and middle-income countries combined. Other regions are increasingly playing larger roles in the growing global smoking epidemic. The WHO Eastern Mediterranean Region (EMRO) now has the highest growth rate in the cigarette market, with more than a one-third increase in cigarette consumption since 2000. Due to its recent dynamic economic development and continued population growth, Africa presents the greatest risk in terms of future growth in tobacco use. Within countries, patterns of cigarette consumption also can vary widely. For example, in many of the countries where few women smoke, smoking rates are often high in males (e.g., in Asia). By contrast, in most developed countries, female smoking rates are typically only a few percentage points below those of males. In many high and middle income countries lower socioeconomic status is a strong predictor of smoking. Smoking rates in the United States have dropped by more than half from 1965 to 2016, falling from 42% to 15.5% of US adults. Australia is cutting their overall smoking consumption faster than most of the developed world, in part due to landmark Plain Packaging Act, which standardized the appearance of cigarette packs. Other countries have considered similar measures. In New Zealand, a bill has been presented to parliament in which the government's associate health minister said "takes away the last means of promoting tobacco as a desirable product." Lights Some cigarettes are marketed as "lights", "milds", or "low-tar". These cigarettes were historically marketed as being less harmful, but there is no research showing that they are any less harmful. The filter design is one of the main differences between light and regular cigarettes, although not all cigarettes contain perforated holes in the filter. In some light cigarettes, the filter is perforated with small holes that theoretically diffuse the tobacco smoke with clean air. In regular cigarettes, the filter does not include these perforations. In ultralight cigarettes, the filter's perforations are larger. The majority of major cigarette manufacturers offer a light, low-tar, or mild cigarette brand. Due to recent U.S. legislation prohibiting the use of these descriptors, tobacco manufacturers are turning to color-coding to allow consumers to differentiate between regular and light brands. Research shows that smoking "light" or "low-tar" cigarettes is just as harmful as smoking other cigarettes. Notable cigarette brands 555 Amber Leaf Army Club Basic Benson & Hedges Barclay Camel Capri Chesterfield Davidoff Dunhill Djarum Dji Sam Soe Doral du Maurier Eclipse Embassy Eve Export Fatima Fortuna Fortune International Gauloises Gitanes Gold Flake Golden Virginia Gold Leaf Kyriazi Freres Kent Kool Lambert and Butler L&M Lark Lucky Strike Marlboro Mayfair Merit Mild Seven More Nat Sherman Natural American Spirit Newport Next Nil Old Gold Pall Mall Parliament Perilly's Peter Stuyvesant Peter Jackson Philip Morris Player's Prince Dunhill Salem Sampoerna Senior Service Smokin Joes Sobranie Sterling Surya Tareyton Vantage Viceroy Virginia Slims West Woodbine Winfield Winston Smoking cessation Smoking cessation (quitting smoking) is the process of discontinuing the practice of tobacco smoking. Quitting can be difficult for many smokers due to the addictive nature of nicotine. The addiction begins when nicotine acts on nicotinic acetylcholine receptors to release neurotransmitters such as dopamine, glutamate, and gamma-aminobutyric acid. Cessation of smoking leads to symptoms of nicotine withdrawal such as anxiety and irritability. Professional smoking cessation support methods generally endeavour to address both nicotine addiction and nicotine withdrawal symptoms. Smoking cessation can be achieved with or without assistance from healthcare professionals or the use of medications. Methods that have been found to be effective include interventions directed at or through health care providers and health care systems; medications including nicotine replacement therapy (NRT) and varenicline; individual and group counselling; and web-based or stand-alone computer programs. Although stopping smoking can cause short-term side effects such as reversible weight gain, smoking cessation services and activities are cost-effective because of the positive health benefits. At the University of Buffalo, researchers found out that fruit and vegetable consumption can help a smoker cut down or even quit smoking A growing number of countries have more ex-smokers than smokers. Early "failure" is a normal part of trying to stop, and more than one attempt at stopping smoking prior to longer-term success is common. NRT, other prescribed pharmaceuticals, and professional counselling or support also help many smokers. However, up to three-quarters of ex-smokers report having quit without assistance ("cold turkey" or cut down then quit), and cessation without professional support or medication may be the most common method used by ex-smokers. The number of nicotinic receptors in the brain returns to the level of a nonsmoker between 6 and 12 weeks after quitting. In 2019, the FDA authorized the selling of low-nicotine cigarettes in hopes of lowering the number of people addicted to nicotine.	Biology and health sciences	Drugs and pharmacology	null
38352	https://en.wikipedia.org/wiki/Rapier	Rapier	A rapier () is a type of sword originally used in Spain (known as -) and Italy (known as spada da lato a striscia). The name designates a sword with a straight, slender and sharply pointed two-edged long blade wielded in one hand. It was widely popular in Western Europe throughout the 16th and 17th centuries as a symbol of nobility or gentleman status. It is called because it was carried as an accessory to clothing, generally used for fashion and as a weapon for dueling, self-defense and as a military side arm. Its name is of Spanish origin and appears recorded for the first time in the Coplas de la panadera, by Juan de Mena, written approximately between 1445 and 1450: As Fencing spread throughout Western Europe, important sources for rapier fencing arose in Spain, known under the term ("dexterity"), in Italy and France. The French small sword or court sword of the 18th century was a direct continuation of this tradition of fencing. Rapier fencing forms part of Historical European Martial Arts. Terminology The origin of the name 'rapier' is Spanish. Its name is a "derisive" description of the Spanish term . The Spanish term refers to a sword used with clothes (, ), due to it being used as an accessory for clothing, usually for fashion and as a self-defense weapon. The English term "rapier" comes from the French and appears both in English and German, near-simultaneously, in the mid-16th century, for a light, long, pointed two-edged sword. It is a loan from Middle French , first recorded in 1474, a nickname meaning . The 16th-century German described what was considered a foreign weapon, imported from Spain, Italy, and France. Du Cange in his Middle Latin dictionary cites a form from a Latin text of 1511. He envisages a derivation . Adelung in his 1798 dictionary records a double meaning for the German verb : on one hand, and on the other. The terms used by the Spanish, Italian and French masters during the heyday of this weapon were simply the equivalent of "sword", i.e. , spada and (). When it was necessary to specify the type of sword the Spanish used . The name was registered for the first time in las Coplas de la panadera, by Juan de Mena, written between 1445 and 1450 approximately. Clements (1997) categorizes thrusting swords with poor cutting abilities as rapiers, and swords with both good thrusting and cutting abilities as cut-and-thrust swords. The term "rapier" is also applied by archaeologists to an unrelated type of Bronze Age sword. Description The word "rapier" generally refers to a relatively long-bladed sword characterized by a protective hilt which is constructed to provide protection for the hand wielding the sword. Some historical rapier samples also feature a broad blade mounted on a typical rapier hilt. The term rapier can be confusing because this hybrid weapon can be categorized as a type of broadsword. While the rapier blade might be broad enough to cut to some degree (but nowhere near that of the wider swords in use around the Middle Ages such as the longsword), it is designed to perform quick and nimble thrusting attacks. The blade might be sharpened along its entire length or sharpened only from the center to the tip (as described by Capoferro). Pallavicini, a rapier master in 1670, strongly advocated using a weapon with two cutting edges. A typical example would weigh and have a relatively long and slender blade of or less in width, or more in length and ending in a sharply pointed tip. The blade length of quite a few historical examples, particularly the Italian rapiers in the early 17th century, is well over and can even reach . The term rapier generally refers to a thrusting sword with a blade longer and thinner than that of the so-called side-sword but much heavier than the small sword, a lighter weapon that would follow in the 18th century and later, but the exact form of the blade and hilt often depends on who is writing and when. It can refer to earlier and the similar espada ropera, through the high rapier period of the 17th century through the small sword and duelling swords; thus context is important in understanding what is meant by the word. (The term side-sword, used among some modern historical martial arts reconstructionists, is a translation from the Italian —a term coined long after the fact by Italian museum curators—and does not refer to the slender, long rapier, but only to the early 16th-century Italian sword with a broader and shorter blade that is considered both its ancestor and contemporary.) Parts of the sword Hilt Rapiers often have complex, sweeping hilts designed to protect the hand wielding the sword. Rings extend forward from the crosspiece. In some later samples, rings are covered with metal plates, eventually evolving into the cup hilts of many later rapiers. There were hardly any samples that featured plates covering the rings prior to the 1600s. Many hilts include a knuckle bow extending down from the crosspiece protecting the grip, which was usually wood wrapped with cord, leather or wire. A large pommel (often decorated) secures the hilt to the weapon and provides some weight to balance the long blade. Blade Various rapier masters divided the blade into two, three, four, five or even nine parts. The forte, strong, is that part of the blade closest to the hilt; in cases where a master divides the blade into an even number of parts, this is the first half of the blade. The debole, weak, is the part of the blade which includes the point and is the second half of the blade when the sword is divided into an even number of parts. However, some rapier masters divided the blade into three parts (or even a multiple of three), in which case the central third of the blade, between the forte and the debole, was often called the medio, mezzo or the terzo. Others used four divisions (Fabris) or even 12 (Thibault). The ricasso is the rear portion of the blade, usually unsharpened. It extends forward from the crosspiece or quillion and then gradually integrates into the thinner and sharper portion of the blade. Overall length There was historical disagreement over how long the ideal rapier should be, with some masters, such as Thibault, denigrating those who recommended longer blades; Thibault's own recommended length was such that the cross of the sword be level with the navel (belly button) when standing naturally with the point resting on the ground. A small number of rapiers with extending blades were made, of which four survive in modern collections. The purpose of the ability is unclear, with suggestions including trying to gain the advantage of surprise in a duel or an attempt to get around laws limiting weapon length. Off-hand weapons Rapiers are single-handed weapons and they were often employed with off-hand bucklers, daggers, cloaks and even second swords to assist with defense. A buckler is a small round shield that was used with other blades as well, such as the arming sword. Capo Ferro's Gran Simulacro depicts use of the weapon with the rotella, which is a significantly bigger shield compared with the buckler. Nevertheless, using the rapier with a parrying dagger is the most common practice, and it has been arguably considered as the most suited and effective accompanying weapon for the rapier. Even though the slender blade of a rapier enables the user to launch a quick attack at a fairly long and advantaged distance between the user and the opponent, and the protective hilt can deflect the opponent's blade, the thrust-oriented weapon is weakened by its bated cutting power and relatively low maneuverability at a closer distance, where the opponent has safely passed the reach of the rapier's deadly point. Therefore, some close-range protection for the user needs to be ensured if the user intends to use the rapier in an optimal way, especially when the opponent uses some slash-oriented sword like a sabre or a broadsword. A parrying dagger not only enables the users to defend in this scenario in which the rapier is not very good at protecting the user, but also enables them to attack in such close distance. History The espada ropera of the 16th century was a cut-and-thrust civilian weapon for self-defense and the duel, while earlier weapons were equally at home on the battlefield. Throughout the 16th century, a variety of new, single-handed civilian weapons were being developed. In 1570, the Italian master Rocco Bonetti first settled in England advocating the use of the rapier for thrusting as opposed to cutting or slashing when engaged in a duel. Nevertheless, the English word "rapier" generally refers to a primarily thrusting weapon, developed by the year 1600 as a result of the geometrical theories of such masters as Camillo Agrippa, Ridolfo Capo Ferro, and Vincentio Saviolo. The rapier became extremely fashionable throughout Europe with the wealthier classes, but was not without its detractors. Some people, such as George Silver, disapproved of its technical potential and the dueling use to which it was put. Allowing for fast reactions, and with a long reach, the rapier was well suited to civilian combat in the 16th and 17th centuries. As military-style cutting and thrusting swords continued to evolve to meet needs on the battlefield, the rapier continued to evolve to meet the needs of civilian combat and decorum, eventually becoming lighter, shorter and less cumbersome to wear. This is when the rapier began to give way to the colichemarde, which was itself later superseded by the small sword which was later superseded by the épée. Noticeably, there were some "war rapiers" that feature a relatively wide blade mounted on a typical rapier hilt during this era. These hybrid swords were used in the military, even on the battlefield. The sword carried by King Gustavus Adolphus in the Thirty Years' War is a typical example of the "war rapier". By the year 1715, the rapier had been largely replaced by the lighter small sword throughout most of Europe, although the former continued to be used, as evidenced by the treatises of Donald McBane (1728), P. J. F. Girard (1736) and Domenico Angelo (1787). The rapier is still used today by officers of the Swiss Guard of the pope. Historical schools of rapier fencing Italy Achille Marozzo, Opera Nova Chiamata Duello, O Vero Fiore dell'Armi de Singulari Abattimenti Offensivi, & Diffensivi1536 Angelo Viggiani dal Montone, Trattato dello Schermo1575 Anonimo Bolognese, L'Arte della Spada (M-345/M-346 Manuscripts)(early or mid 16th century) date it to "about 1550" Antonio Manciolino, Opera Nova per Imparare a Combattere, & Schermire d'ogni sorte Armi1531 Bondi di Mazo, La Spada Maestra1696 Camillo Agrippa, Trattato di Scientia d'Arme con un Dialogo di Filosofia1553 Francesco Alfieri, La Scherma di Francesco Alfieri1640 Francesco Antonio Marcelli, Regole della Scherma1686 Giacomo di Grassi, Ragion di Adoprar Sicuramente l'Arme si da Offesa, come da Difesa1570 Giovanni dall'Agocchie, Dell'Arte di Scrimia1572 Giuseppe Morsicato Pallavicini, La Scherma Illustrata1670 Marco Docciolini, Trattato in Materia di Scherma1601 Nicoletto Giganti, Scola overo Teatro1606 Ridolfo Capo Ferro, Gran Simulacro dell'Arte e dell'Uso della Scherma1610 Salvator Fabris, De lo Schermo ovvero Scienza d'Armi1606 Spain Jerónimo Sánchez de Carranza, De la Filosofía de las Armas (1569) Luis Pacheco de Narváez, Libro de las Grandezas de la Espada (1599) The Netherlands Girard Thibault, Academie de l'Espée (1630) France André Desbordes, Discours de la théorie et de la pratique de l'excellence des armes (1610) Charles Besnard, Le maistre d'arme liberal (1653) François Dancie, Discours des armes et methode pour bien tirer de l'espée et poignard () and L'Espee de combat (1623) England Joseph Swetnam, The Schoole of the Noble and Worthy Science of Defence (1617) The Pallas Armata (1639) Vincentio Saviolo, His Practise 1595 Germany Jakob Sutor, Künstliches Fechtbuch (1612) Joachim Meyer, Thorough Descriptions of the free Knightly and Noble Art of Fencing (1570) Johannes Georgius Bruchius (1671) Paulus Hector Mair, Opus Amplissimum de Arte Athletica (1542) The classical fencing tradition Classical fencing schools claim to have inherited aspects of rapier forms in their systems. In 1885, fencing scholar Egerton Castle wrote "there is little doubt that the French system of fencing can be traced, at its origin, to the ancient Italian swordsmanship; the modern Italian school being of course derived in an uninterrupted manner from the same source." Castle went on to note that "the Italians have preserved the rapier form, with cup, pas d'ane, and quillons, but with a slender quadrangular blade."	Technology	Melee weapons	null
38390	https://en.wikipedia.org/wiki/Dementia	Dementia	Dementia is a syndrome associated with many neurodegenerative diseases, characterized by a general decline in cognitive abilities that affects a person's ability to perform everyday activities. This typically involves problems with memory, thinking, behavior, and motor control. Aside from memory impairment and a disruption in thought patterns, the most common symptoms of dementia include emotional problems, difficulties with language, and decreased motivation. The symptoms may be described as occurring in a continuum over several stages. Dementia ultimately has a significant effect on the individual, their caregivers, and their social relationships in general. A diagnosis of dementia requires the observation of a change from a person's usual mental functioning and a greater cognitive decline than might be caused by the normal aging process. Several diseases and injuries to the brain, such as a stroke, can give rise to dementia. However, the most common cause is Alzheimer's disease, a neurodegenerative disorder. The Diagnostic and Statistical Manual of Mental Disorders, Fifth Edition (DSM-5), has re-described dementia as a mild or major neurocognitive disorder with varying degrees of severity and many causative subtypes. The International Classification of Diseases (ICD-11) also classifies dementia as a neurocognitive disorder (NCD) with many forms or subclasses. Dementia is listed as an acquired brain syndrome, marked by a decline in cognitive function, and is contrasted with neurodevelopmental disorders. It is also described as a spectrum of disorders with causative subtypes of dementia based on a known disorder, such as Parkinson's disease for Parkinson's disease dementia, Huntington's disease for Huntington's disease dementia, vascular disease for vascular dementia, HIV infection causing HIV dementia, frontotemporal lobar degeneration for frontotemporal dementia, Lewy body disease for dementia with Lewy bodies, and prion diseases. Subtypes of neurodegenerative dementias may also be based on the underlying pathology of misfolded proteins, such as synucleinopathies and tauopathies. The coexistence of more than one type of dementia is known as mixed dementia. Many neurocognitive disorders may be caused by another medical condition or disorder, including brain tumours and subdural hematoma, endocrine disorders such as hypothyroidism and hypoglycemia, nutritional deficiencies including thiamine and niacin, infections, immune disorders, liver or kidney failure, metabolic disorders such as Kufs disease, some leukodystrophies, and neurological disorders such as epilepsy and multiple sclerosis. Some of the neurocognitive deficits may sometimes show improvement with treatment of the causative medical condition. Diagnosis of dementia is usually based on history of the illness and cognitive testing with imaging. Blood tests may be taken to rule out other possible causes that may be reversible, such as hypothyroidism (an underactive thyroid), and to determine the dementia subtype. One commonly used cognitive test is the mini–mental state examination. Although the greatest risk factor for developing dementia is aging, dementia is not a normal part of the aging process; many people aged 90 and above show no signs of dementia. Several risk factors for dementia, such as smoking and obesity, are preventable by lifestyle changes. Screening the general older population for the disorder is not seen to affect the outcome. Dementia is currently the seventh leading cause of death worldwide and has 10 million new cases reported every year (approximately one every three seconds). There is no known cure for dementia. Acetylcholinesterase inhibitors such as donepezil are often used and may be beneficial in mild to moderate disorder, but the overall benefit may be minor. There are many measures that can improve the quality of life of a person with dementia and their caregivers. Cognitive and behavioral interventions may be appropriate for treating the associated symptoms of depression. Signs and symptoms The signs and symptoms of dementia are termed as the neuropsychiatric symptoms—also known as the behavioral and psychological symptoms—of dementia. The behavioral symptoms can include agitation, restlessness, inappropriate behavior, sexual disinhibition, and verbal or physical aggression. These symptoms may result from impairments in cognitive inhibition. The psychological symptoms can include depression, hallucinations (most often visual), delusions, apathy, and anxiety. The most commonly affected areas of brain function include memory, language, attention, problem solving, and visuospatial function affecting perception and orientation. The symptoms progress at a continuous rate over several stages, and they vary across the dementia subtypes. Most types of dementia are slowly progressive with some deterioration of the brain well established before signs of the disorder become apparent. There are often other conditions present, such as high blood pressure or diabetes, and there can sometimes be as many as four of these comorbidities. Signs of dementia include getting lost in a familiar neighborhood, using unusual words to refer to familiar objects, forgetting the name of a close family member or friend, forgetting old memories, and being unable to complete tasks independently. People with developing dementia often fall behind on bill payments; specifically mortgage and credit cards, and a crashing credit score can be an early indicator of the disease. People with dementia are more likely to have problems with incontinence than those of a comparable age without dementia; they are three times more likely to have urinary incontinence and four times more likely to have fecal incontinence. Stages The course of dementia is often described in four stages – pre-dementia, early, middle, and late, that show a pattern of progressive cognitive and functional impairment. More detailed descriptions can be arrived at by the use of numeric scales. These scales include the Global Deterioration Scale (GDS or Reisberg Scale), the Functional Assessment Staging Tool (FAST), and the Clinical Dementia Rating (CDR). Using the GDS, which more accurately identifies each stage of the disease progression, a more detailed course is described in seven stages – two of which are broken down further into five and six degrees. Stage 7(f) is the final stage. Pre-dementia Pre-dementia includes pre-clinical and prodromal stages. The latter stage includes mild cognitive impairment (MCI), delirium-onset, and psychiatric-onset presentations. Pre-clinical Sensory dysfunction is claimed for the pre-clinical stage, which may precede the first clinical signs of dementia by up to ten years. Most notably the sense of smell is lost, associated with depression and a loss of appetite leading to poor nutrition. It is suggested that this dysfunction may come about because the olfactory epithelium is exposed to the environment, and the lack of blood–brain barrier protection allows toxic elements to enter and cause damage to the chemosensory networks. Prodromal Pre-dementia states considered as prodromal are mild cognitive impairment (MCI) and mild behavioral impairment (MBI). Signs and symptoms at the prodromal stage may be subtle, and the early signs often become apparent only in hindsight. Of those diagnosed with MCI, 70% later progress to dementia. In mild cognitive impairment, changes in the person's brain have been happening for a long time, but the symptoms are just beginning to appear. These problems, however, are not severe enough to affect daily function. If and when they do, the diagnosis becomes dementia. The person may have some memory problems and trouble finding words, but they can solve everyday problems and competently handle their life affairs. During this stage, it is ideal to ensure that advance care planning has occurred to protect the person's wishes. Advance directives exist that are specific to sufferers of dementia; these can be particularly helpful in addressing the decisions related to feeding which come with the progression of the illness. Mild cognitive impairment has been relisted in both DSM-5 and ICD-11 as "mild neurocognitive disorders", i.e. milder forms of the major neurocognitive disorder (dementia) subtypes. Kynurenine is a metabolite of tryptophan that regulates microbiome signaling, immune cell response, and neuronal excitation. A disruption in the kynurenine pathway may be associated with the neuropsychiatric symptoms and cognitive prognosis in mild dementia. Early In the early stage of dementia, symptoms become noticeable to other people. In addition, the symptoms begin to interfere with daily activities, and will register a score on a mini–mental state examination (MMSE). MMSE scores are set at 24 to 30 for a normal cognitive rating and lower scores reflect severity of symptoms. The symptoms are dependent on the type of dementia. More complicated chores and tasks around the house or at work become more difficult. The person can usually still take care of themselves but may forget things like taking pills or doing laundry and may need prompting or reminders. The symptoms of early dementia usually include memory difficulty, but can also include some word-finding problems, and problems with executive functions of planning and organization. Managing finances may prove difficult. Other signs might be getting lost in new places, repeating things, and personality changes. In some types of dementia, such as dementia with Lewy bodies and frontotemporal dementia, personality changes and difficulty with organization and planning may be the first signs. Middle As dementia progresses, initial symptoms generally worsen. The rate of decline is different for each person. MMSE scores between 6 and 17 signal moderate dementia. For example, people with moderate Alzheimer's dementia lose almost all new information. People with dementia may be severely impaired in solving problems, and their social judgment is often impaired. They cannot usually function outside their own home, and generally should not be left alone. They may be able to do simple chores around the house but not much else, and begin to require assistance for personal care and hygiene beyond simple reminders. A lack of insight into having the condition will become evident. Late People with late-stage dementia typically turn increasingly inward and need assistance with most or all of their personal care. People with dementia in the late stages usually need 24-hour supervision to ensure their personal safety, and meeting of basic needs. If left unsupervised, they may wander or fall; may not recognize common dangers such as a hot stove; or may not realize that they need to use the bathroom and become incontinent. They may not want to get out of bed, or may need assistance doing so. Commonly, the person no longer recognizes familiar faces. They may have significant changes in sleeping habits or have trouble sleeping at all. Changes in eating frequently occur. Cognitive awareness is needed for eating and swallowing and progressive cognitive decline results in eating and swallowing difficulties. This can cause food to be refused, or choked on, and help with feeding will often be required. For ease of feeding, food may be liquidized into a thick purée. They may also struggle to walk, particularly among those with Alzheimer's disease. In some cases, terminal lucidity, a form of paradoxical lucidity, occurs immediately before death; in this phenomenon, there is an unexpected recovery of mental clarity. Causes Many causes of dementia are neurodegenerative, and protein misfolding is a cardinal feature of these. Other common causes include vascular dementia, dementia with Lewy bodies, frontotemporal dementia, and mixed dementia (commonly Alzheimer's disease and vascular dementia). Less common causes include normal pressure hydrocephalus, Parkinson's disease dementia, syphilis, HIV, and Creutzfeldt–Jakob disease. Alzheimer's disease Alzheimer's disease accounts for 60–70% of cases of dementia worldwide. The most common symptoms of Alzheimer's disease are short-term memory loss and word-finding difficulties. Trouble with visuospatial functioning (getting lost often), reasoning, judgment and insight fail. Insight refers to whether or not the person realizes they have memory problems. The part of the brain most affected by Alzheimer's is the hippocampus. Other parts that show atrophy (shrinking) include the temporal and parietal lobes. Although this pattern of brain shrinkage suggests Alzheimer's, it is variable and a brain scan is insufficient for a diagnosis. Little is known about the events that occur during and that actually cause Alzheimer's disease. This is due to the fact that, historically, brain tissue from patients with the disease could only be studied after the person's death. Brain scans can now help diagnose and distinguish between different kinds of dementia and show severity. These include magnetic resonance imaging (MRI), computerized tomography (CT), and positron emission tomography (PET). However, it is known that one of the first aspects of Alzheimer's disease is overproduction of amyloid. Extracellular senile plaques (SPs), consisting of beta-amyloid (Aβ) peptides, and intracellular neurofibrillary tangles (NFTs) that are formed by hyperphosphorylated tau proteins, are two well-established pathological hallmarks of AD. Amyloid causes inflammation around the senile plaques of the brain, and too much buildup of this inflammation leads to changes in the brain that cannot be controlled, leading to the symptoms of Alzheimer's. Several articles have been published on a possible relationship (as an either primary cause or exacerbation of Alzheimer's disease) between general anesthesia and Alzheimer's in specifically the elderly. Vascular Vascular dementia accounts for at least 20% of dementia cases, making it the second most common type. It is caused by disease or injury affecting the blood supply to the brain, typically involving a series of mini-strokes. The symptoms of this dementia depend on where in the brain the strokes occurred and whether the blood vessels affected were large or small. Repeated injury can cause progressive dementia over time, while a single injury located in an area critical for cognition such as the hippocampus, or thalamus, can lead to sudden cognitive decline. Elements of vascular dementia may be present in all other forms of dementia. Brain scans may show evidence of multiple strokes of different sizes in various locations. People with vascular dementia tend to have risk factors for disease of the blood vessels, such as tobacco use, high blood pressure, atrial fibrillation, high cholesterol, diabetes, or other signs of vascular disease such as a previous heart attack or angina. Lewy bodies The prodromal symptoms of dementia with Lewy bodies (DLB) include mild cognitive impairment, and delirium onset. The symptoms of DLB are more frequent, more severe, and earlier presenting than in the other dementia subtypes. Dementia with Lewy bodies has the primary symptoms of fluctuating cognition, alertness or attention; REM sleep behavior disorder (RBD); one or more of the main features of parkinsonism, not due to medication or stroke; and repeated visual hallucinations. The visual hallucinations in DLB are generally vivid hallucinations of people or animals and they often occur when someone is about to fall asleep or wake up. Other prominent symptoms include problems with planning (executive function) and difficulty with visual-spatial function, and disruption in autonomic bodily functions. Abnormal sleep behaviors may begin before cognitive decline is observed and are a core feature of DLB. RBD is diagnosed either by sleep study recording or, when sleep studies cannot be performed, by medical history and validated questionnaires. Parkinson's disease Parkinson's disease is associated with Lewy body dementia that often progresses to Parkinson's disease dementia following a period of dementia-free Parkinson's disease. Frontotemporal Frontotemporal dementias (FTDs) are characterized by drastic personality changes and language difficulties. In all FTDs, the person has a relatively early social withdrawal and early lack of insight. Memory problems are not a main feature. There are six main types of FTD. The first has major symptoms in personality and behavior. This is called behavioral variant FTD (bv-FTD) and is the most common. The hallmark feature of bv-FTD is impulsive behavior, and this can be detected in pre-dementia states. In bv-FTD, the person shows a change in personal hygiene, becomes rigid in their thinking, and rarely acknowledges problems; they are socially withdrawn, and often have a drastic increase in appetite. They may become socially inappropriate. For example, they may make inappropriate sexual comments, or may begin using pornography openly. One of the most common signs is apathy, or not caring about anything. Apathy, however, is a common symptom in many dementias. Two types of FTD feature aphasia (language problems) as the main symptom. One type is called semantic variant primary progressive aphasia (SV-PPA). The main feature of this is the loss of the meaning of words. It may begin with difficulty naming things. The person eventually may lose the meaning of objects as well. For example, a drawing of a bird, dog, and an airplane in someone with FTD may all appear almost the same. In a classic test for this, a patient is shown a picture of a pyramid and below it a picture of both a palm tree and a pine tree. The person is asked to say which one goes best with the pyramid. In SV-PPA the person cannot answer that question. The other type is called non-fluent agrammatic variant primary progressive aphasia (NFA-PPA). This is mainly a problem with producing speech. They have trouble finding the right words, but mostly they have a difficulty coordinating the muscles they need to speak. Eventually, someone with NFA-PPA only uses one-syllable words or may become totally mute. A frontotemporal dementia associated with amyotrophic lateral sclerosis (ALS) known as (FTD-ALS) includes the symptoms of FTD (behavior, language and movement problems) co-occurring with amyotrophic lateral sclerosis (loss of motor neurons). Two FTD-related disorders are progressive supranuclear palsy (also classed as a Parkinson-plus syndrome), and corticobasal degeneration. These disorders are tau-associated. Huntington's disease Huntington's disease is a neurodegenerative disease caused by mutations in a single gene HTT, that encodes for huntingtin protein. Symptoms include cognitive impairment and this usually declines further into dementia. The first main symptoms of Huntington's disease often include: difficulty concentrating memory lapses depression - this can include low mood, lack of interest in things, or just abnormal feelings of hopelessness stumbling and clumsiness that is out of the ordinary mood swings, such as irritability or aggressive behavior to insignificant things HIV HIV-associated dementia results as a late stage from HIV infection, and mostly affects younger people. The essential features of HIV-associated dementia are disabling cognitive impairment accompanied by motor dysfunction, speech problems and behavioral change. Cognitive impairment is characterised by mental slowness, trouble with memory and poor concentration. Motor symptoms include a loss of fine motor control leading to clumsiness, poor balance and tremors. Behavioral changes may include apathy, lethargy and diminished emotional responses and spontaneity. Histopathologically, it is identified by the infiltration of monocytes and macrophages into the central nervous system (CNS), gliosis, pallor of myelin sheaths, abnormalities of dendritic processes and neuronal loss. Creutzfeldt–Jakob disease Creutzfeldt–Jakob disease is a rapidly progressive prion disease that typically causes dementia that worsens over weeks to months. Prions are disease-causing pathogens created from abnormal proteins. Alcoholism Alcohol-related dementia, also called alcohol-related brain damage, occurs as a result of excessive use of alcohol particularly as a substance abuse disorder. Different factors can be involved in this development including thiamine deficiency and age vulnerability. A degree of brain damage is seen in more than 70% of those with alcohol use disorder. Brain regions affected are similar to those that are affected by aging, and also by Alzheimer's disease. Regions showing loss of volume include the frontal, temporal, and parietal lobes, as well as the cerebellum, thalamus, and hippocampus. This loss can be more notable, with greater cognitive impairments seen in those aged 65 years and older. Mixed dementia More than one type of dementia, known as mixed dementia, may exist together in about 10% of dementia cases. The most common type of mixed dementia is Alzheimer's disease and vascular dementia. This particular type of mixed dementia's main onsets are a mixture of old age, high blood pressure, and damage to blood vessels in the brain. Diagnosis of mixed dementia can be difficult, as often only one type will predominate. This makes the treatment of people with mixed dementia uncommon, with many people missing out on potentially helpful treatments. Mixed dementia can mean that symptoms onset earlier, and worsen more quickly since more parts of the brain will be affected. Other Chronic inflammatory conditions that may affect the brain and cognition include Behçet's disease, multiple sclerosis, sarcoidosis, Sjögren's syndrome, lupus, celiac disease, and non-celiac gluten sensitivity. These types of dementias can rapidly progress, but usually have a good response to early treatment. This consists of immunomodulators or steroid administration, or in certain cases, the elimination of the causative agent. Celiac disease does not seem to raise the risk of dementia in general but it may increase the risk of vascular dementia. Both celiac disease or non-celiac gluten sensitivity might raise the risk of cognitive impairment which can be one of the early signs of subsequent dementia. A strict gluten-free diet started early may protect against dementia associated with gluten-related disorders. Cases of easily reversible dementia include hypothyroidism, vitamin B12 deficiency, Lyme disease, and neurosyphilis. For Lyme disease and neurosyphilis, testing should be done if risk factors are present. Because risk factors are often difficult to determine, testing for neurosyphilis and Lyme disease, as well as other mentioned factors, may be undertaken as a matter of course where dementia is suspected. Many other medical and neurological conditions include dementia only late in the illness. For example, a proportion of patients with Parkinson's disease develop dementia, though widely varying figures are quoted for this proportion. When dementia occurs in Parkinson's disease, the underlying cause may be dementia with Lewy bodies or Alzheimer's disease, or both. Cognitive impairment also occurs in the Parkinson-plus syndromes of progressive supranuclear palsy and corticobasal degeneration (and the same underlying pathology may cause the clinical syndromes of frontotemporal lobar degeneration). Although the acute porphyrias may cause episodes of confusion and psychiatric disturbance, dementia is a rare feature of these rare diseases. Limbic-predominant age-related TDP-43 encephalopathy (LATE) is a type of dementia that primarily affects people in their 80s or 90s and in which TDP-43 protein deposits in the limbic portion of the brain. Hereditary disorders that can also cause dementia include: some metabolic disorders such as lysosomal storage disorders, leukodystrophies, and spinocerebellar ataxias. Persistent loneliness may significantly increase the risk of dementia. Loneliness is associated with a 31% higher likelihood of developing any form of dementia, and can also raise the risk of cognitive impairment by 15%. Diagnosis Symptoms are similar across dementia types and it is difficult to diagnose by symptoms alone. Diagnosis may be aided by brain scanning techniques. In many cases, the diagnosis requires a brain biopsy to become final, but this is rarely recommended (though it can be performed at autopsy). In those who are getting older, general screening for cognitive impairment using cognitive testing or early diagnosis of dementia has not been shown to improve outcomes. However, screening exams are useful in 65+ persons with memory complaints. Normally, symptoms must be present for at least six months to support a diagnosis. Cognitive dysfunction of shorter duration is called delirium. Delirium can be easily confused with dementia due to similar symptoms. Delirium is characterized by a sudden onset, fluctuating course, a short duration (often lasting from hours to weeks), and is primarily related to a somatic (or medical) disturbance. In comparison, dementia has typically a long, slow onset (except in the cases of a stroke or trauma), slow decline of mental functioning, as well as a longer trajectory (from months to years). Some mental illnesses, including depression and psychosis, may produce symptoms that must be differentiated from both delirium and dementia. These are differently diagnosed as pseudodementias, and any dementia evaluation needs to include a depression screening such as the Neuropsychiatric Inventory or the Geriatric Depression Scale. Physicians used to think that people with memory complaints had depression and not dementia (because they thought that those with dementia are generally unaware of their memory problems). However, researchers have realized that many older people with memory complaints in fact have mild cognitive impairment the earliest stage of dementia. Depression should always remain high on the list of possibilities, however, for an elderly person with memory trouble. Changes in thinking, hearing and vision are associated with normal ageing and can cause problems when diagnosing dementia due to the similarities. Given the challenging nature of predicting the onset of dementia and making a dementia diagnosis clinical decision making aids underpinned by machine learning and artificial intelligence have the potential to enhance clinical practice. Cognitive testing Various brief cognitive tests (5–15 minutes) have reasonable reliability to screen for dementia, but may be affected by factors such as age, education and ethnicity. Age and education have a significant influence on the diagnosis of dementia. For example, Individuals with lower education are more likely to be diagnosed with dementia than their educated counterparts. While many tests have been studied, presently the mini mental state examination (MMSE) is the best studied and most commonly used. The MMSE is a useful tool for helping to diagnose dementia if the results are interpreted along with an assessment of a person's personality, their ability to perform activities of daily living, and their behaviour. Other cognitive tests include the abbreviated mental test score (AMTS), the, "modified mini–mental state examination" (3MS), the Cognitive Abilities Screening Instrument (CASI), the Trail-making test, and the clock drawing test. The MoCA (Montreal Cognitive Assessment) is a reliable screening test and is available online for free in 35 different languages. The MoCA has also been shown somewhat better at detecting mild cognitive impairment than the MMSE. People with hearing loss, which commonly occurs alongside dementia, score worse in the MoCA test, which could lead to a false diagnosis of dementia. Researchers have developed an adapted version of the MoCA test, which is accurate and reliable and avoids the need for people to listen and respond to questions. The AD-8 – a screening questionnaire used to assess changes in function related to cognitive decline – is potentially useful, but is not diagnostic, is variable, and has risk of bias. An integrated cognitive assessment (CognICA) is a five-minute test that is highly sensitive to the early stages of dementia, and uses an application deliverable to an iPad. Previously in use in the UK, in 2021 CognICA was given FDA approval for its commercial use as a medical device. Another approach to screening for dementia is to ask an informant (relative or other supporter) to fill out a questionnaire about the person's everyday cognitive functioning. Informant questionnaires provide complementary information to brief cognitive tests. Probably the best known questionnaire of this sort is the Informant Questionnaire on Cognitive Decline in the Elderly (IQCODE). Evidence is insufficient to determine how accurate the IQCODE is for diagnosing or predicting dementia. The Alzheimer's Disease Caregiver Questionnaire is another tool. It is about 90% accurate for Alzheimer's when by a caregiver. The General Practitioner Assessment Of Cognition combines both a patient assessment and an informant interview. It was specifically designed for use in the primary care setting. Clinical neuropsychologists provide diagnostic consultation following administration of a full battery of cognitive testing, often lasting several hours, to determine functional patterns of decline associated with varying types of dementia. Tests of memory, executive function, processing speed, attention and language skills are relevant, as well as tests of emotional and psychological adjustment. These tests assist with ruling out other etiologies and determining relative cognitive decline over time or from estimates of prior cognitive abilities. Laboratory tests Routine blood tests are usually performed to rule out treatable causes. These include tests for vitamin B12, folic acid, thyroid-stimulating hormone (TSH), C-reactive protein, full blood count, electrolytes, calcium, renal function, and liver enzymes. Abnormalities may suggest vitamin deficiency, infection, or other problems that commonly cause confusion or disorientation in the elderly. Imaging A CT scan or MRI scan is commonly performed to possibly find either normal pressure hydrocephalus, a potentially reversible cause of dementia, or connected tumor. The scans can also yield information relevant to other types of dementia, such as infarction (stroke) that would point at a vascular type of dementia. These tests do not pick up diffuse metabolic changes associated with dementia in a person who shows no gross neurological problems (such as paralysis or weakness) on a neurological exam. The functional neuroimaging modalities of SPECT and PET are more useful in assessing long-standing cognitive dysfunction, since they have shown similar ability to diagnose dementia as a clinical exam and cognitive testing. The ability of SPECT to differentiate vascular dementia from Alzheimer's disease, appears superior to differentiation by clinical exam. The value of PiB-PET imaging using Pittsburgh compound B (PiB) as a radiotracer has been established in predictive diagnosis, particularly Alzheimer's disease. Prevention Risk factors Risk factors for dementia include high blood pressure, high levels of LDL cholesterol, vision loss, hearing loss, smoking, obesity, depression, inactivity, diabetes, lower levels of education and low social contact. Over-indulgence in alcohol, lack of sleep, anemia, traumatic brain injury, and air pollution can also increase the chance of developing dementia. Many of these risk factors, including the lower level of education, smoking, physical inactivity and diabetes, are modifiable. Several of the group are known as vascular risk factors that may be possible to be reduced or eliminated. Managing these risk factors can reduce the risk of dementia in individuals in their late midlife or older age. A reduction in a number of these risk factors can give a positive outcome. The decreased risk achieved by adopting a healthy lifestyle is seen even in those with a high genetic risk. In addition to the above risk factors, other psychological features, including certain personality traits (high neuroticism, and low conscientiousness), low purpose in life, and high loneliness, are risk factors for Alzheimer's disease and related dementias. For example, based on the English Longitudinal Study of Ageing (ELSA), research found that loneliness in older people can increase the risk of dementia by one-third. Not having a partner (being single, divorced, or widowed) can double the risk of dementia. However, having two or three closer relationships might reduce the risk by three-fifths. The two most modifiable risk factors for dementia are physical inactivity and lack of cognitive stimulation. Physical activity, in particular aerobic exercise, is associated with a reduction in age-related brain tissue loss, and neurotoxic factors thereby preserving brain volume and neuronal integrity. Cognitive activity strengthens neural plasticity and together they help to support cognitive reserve. The neglect of these risk factors diminishes this reserve. Sensory impairments of vision and hearing are modifiable risk factors for dementia. These impairments may precede the cognitive symptoms of Alzheimer's disease for example, by many years. Hearing loss may lead to social isolation which negatively affects cognition. Social isolation is also identified as a modifiable risk factor. Age-related hearing loss in midlife is linked to cognitive impairment in late life, and is seen as a risk factor for the development of Alzheimer's disease and dementia. Such hearing loss may be caused by a central auditory processing disorder that makes the understanding of speech against background noise difficult. Age-related hearing loss is characterised by slowed central processing of auditory information. Worldwide, mid-life hearing loss may account for around 9% of dementia cases. Frailty may increase the risk of cognitive decline, and dementia, and the inverse also holds of cognitive impairment increasing the risk of frailty. Prevention of frailty may help to prevent cognitive decline. There are no medications that can prevent cognitive decline and dementia. However blood pressure lowering medications might decrease the risk of dementia or cognitive problems by around 0.5%. Economic disadvantage has been shown to have a strong link to higher dementia prevalence, which cannot yet be fully explained by other risk factors. Dental health Limited evidence links poor oral health to cognitive decline. However, failure to perform tooth brushing and gingival inflammation can be used as dementia risk predictors. Oral bacteria The link between Alzheimer's and gum disease is oral bacteria. In the oral cavity, bacterial species include P. gingivalis, F. nucleatum, P. intermedia, and T. forsythia. Six oral treponema spirochetes have been examined in the brains of Alzheimer's patients. Spirochetes are neurotropic in nature, meaning they act to destroy nerve tissue and create inflammation. Inflammatory pathogens are an indicator of Alzheimer's disease and bacteria related to gum disease have been found in the brains of patients with Alzheimer's disease. The bacteria invade nerve tissue in the brain, increasing the permeability of the blood–brain barrier and promoting the onset of Alzheimer's. Individuals with a plethora of tooth plaque risk cognitive decline. Poor oral hygiene can have an adverse effect on speech and nutrition, causing general and cognitive health decline. Oral viruses Herpes simplex virus (HSV) has been found in more than 70% of those aged over 50. HSV persists in the peripheral nervous system and can be triggered by stress, illness or fatigue. High proportions of viral-associated proteins in amyloid plaques or neurofibrillary tangles (NFTs) confirm the involvement of HSV-1 in Alzheimer's disease pathology. NFTs are known as the primary marker of Alzheimer's disease. HSV-1 produces the main components of NFTs. Diet Diet is seen to be a modifiable risk factor for the development of dementia. Thiamine deficiency is identified to increase the risk of Alzheimer's disease in adults. The role of thiamine in brain physiology is unique and essential for the normal cognitive function of older people. Many dietary choices of the elderly population, including the higher intake of gluten-free products, compromise the intake of thiamine as these products are not fortified with thiamine. The Mediterranean and DASH diets are both associated with less cognitive decline. A different approach has been to incorporate elements of both of these diets into one known as the MIND diet. These diets are generally low in saturated fats while providing a good source of carbohydrates, mainly those that help stabilize blood sugar and insulin levels. Raised blood sugar levels over a long time, can damage nerves and cause memory problems if they are not managed. Nutritional factors associated with the proposed diets for reducing dementia risk include unsaturated fatty acids, vitamin E, vitamin C, flavonoids, vitamin B, and vitamin D. A study conducted at the University of Exeter in the United Kingdom seems to have confirmed these findings with fruits, vegetables, whole grains, and healthy fats creating an optimum diet that can help reduce the risk of dementia by roughly 25%. The MIND diet may be more protective but further studies are needed. The Mediterranean diet seems to be more protective against Alzheimer's than DASH but there are no consistent findings against dementia in general. The role of olive oil needs further study as it may be one of the most important components in reducing the risk of cognitive decline and dementia. In those with celiac disease or non-celiac gluten sensitivity, a strict gluten-free diet may relieve the symptoms given a mild cognitive impairment. Once dementia is advanced no evidence suggests that a gluten-free diet is useful. Omega-3 fatty acid supplements do not appear to benefit or harm people with mild to moderate symptoms. However, there is good evidence that omega-3 incorporation into the diet is of benefit in treating depression, a common symptom, and potentially modifiable risk factor for dementia. Management There are limited options for treating dementia, with most approaches focused on managing or reducing individual symptoms. There are no treatment options available to delay the onset of dementia. Acetylcholinesterase inhibitors are often used early in the disorder course; however, benefit is generally small. More than half of people with dementia may experience psychological or behavioral symptoms including agitation, sleep problems, aggression, and/or psychosis. Treatment for these symptoms is aimed at reducing the person's distress and keeping the person safe. Treatments other than medication appear to be better for agitation and aggression. Cognitive and behavioral interventions may be appropriate. Some evidence suggests that education and support for the person with dementia, as well as caregivers and family members, improves outcomes. Palliative care interventions may lead to improvements in comfort in dying, but the evidence is low. Exercise programs are beneficial with respect to activities of daily living, and potentially improve dementia. The effect of therapies can be evaluated for example by assessing agitation using the Cohen-Mansfield Agitation Inventory (CMAI); by assessing mood and engagement with the Menorah Park Engagement Scale (MPES); and the Observed Emotion Rating Scale (OERS) or by assessing indicators for depression using the Cornell Scale for Depression in Dementia (CSDD) or a simplified version thereof. Often overlooked in treating and managing dementia is the role of the caregiver and what is known about how they can support multiple interventions. Findings from a 2021 systematic review of the literature found caregivers of people with dementia in nursing homes do not have sufficient tools or clinical guidance for behavioral and psychological symptoms of dementia (BPSD) along with medication use. Simple measures like talking to people about their interests can improve the quality of life for care home residents living with dementia. A programme showed that such simple measures reduced residents' agitation and depression. They also needed fewer GP visits and hospital admissions, which also meant that the programme was cost-saving. Psychological and psychosocial therapies Psychological therapies for dementia include some limited evidence for reminiscence therapy (namely, some positive effects in the areas of quality of life, cognition, communication and mood – the first three particularly in care home settings), some benefit for cognitive reframing for caretakers, unclear evidence for validation therapy and tentative evidence for mental exercises, such as cognitive stimulation programs for people with mild to moderate dementia. Offering personally tailored activities may help reduce challenging behavior and may improve quality of life. It is not clear if personally tailored activities have an impact on affect or improve for the quality of life for the caregiver. Adult daycare centers as well as special care units in nursing homes often provide specialized care for dementia patients. Daycare centers offer supervision, recreation, meals, and limited health care to participants, as well as providing respite for caregivers. In addition, home care can provide one-to-one support and care in the home allowing for more individualized attention that is needed as the disorder progresses. Psychiatric nurses can make a distinctive contribution to people's mental health. Since dementia impairs normal communication due to changes in receptive and expressive language, as well as the ability to plan and problem solve, agitated behavior is often a form of communication for the person with dementia. Actively searching for a potential cause, such as pain, physical illness, or overstimulation can be helpful in reducing agitation. Additionally, using an "ABC analysis of behavior" can be a useful tool for understanding behavior in people with dementia. It involves looking at the antecedents (A), behavior (B), and consequences (C) associated with an event to help define the problem and prevent further incidents that may arise if the person's needs are misunderstood. The strongest evidence for non-pharmacological therapies for the management of changed behaviors in dementia is for using such approaches. Low quality evidence suggests that regular (at least five sessions of) music therapy may help institutionalized residents. It may reduce depressive symptoms and improve overall behaviors. It may also supply a beneficial effect on emotional well-being and quality of life, as well as reduce anxiety. In 2003, The Alzheimer's Society established 'Singing for the Brain' (SftB) a project based on pilot studies which suggested that the activity encouraged participation and facilitated the learning of new songs. The sessions combine aspects of reminiscence therapy and music. Musical and interpersonal connectedness can underscore the value of the person and improve quality of life. Some London hospitals found that using color, designs, pictures and lights helped people with dementia adjust to being at the hospital. These adjustments to the layout of the dementia wings at these hospitals helped patients by preventing confusion. Life story work as part of reminiscence therapy, and video biographies have been found to address the needs of clients and their caregivers in various ways, offering the client the opportunity to leave a legacy and enhance their personhood and also benefitting youth who participate in such work. Such interventions can be more beneficial when undertaken at a relatively early stage of dementia. They may also be problematic in those who have difficulties in processing past experiences Animal-assisted therapy has been found to be helpful. Drawbacks may be that pets are not always welcomed in a communal space in the care setting. An animal may pose a risk to residents, or may be perceived to be dangerous. Certain animals may also be regarded as "unclean" or "dangerous" by some cultural groups. Occupational therapy also addresses psychological and psychosocial needs of patients with dementia through improving daily occupational performance and caregivers' competence. When compensatory intervention strategies are added to their daily routine, the level of performance is enhanced and reduces the burden commonly placed on their caregivers. Occupational therapists can also work with other disciplines to create a client centered intervention. To manage cognitive disability, and coping with behavioral and psychological symptoms of dementia, combined occupational and behavioral therapies can support patients with dementia even further. Cognitive training and rehabilitation There is no strong evidence to suggest that cognitive training is beneficial for people with Parkinson's disease, dementia, or mild cognitive impairment. However, a 2023 review found that cognitive rehabilitation may be effective in helping individuals with mild to moderate dementia to manage their daily activities. Personally tailored activities Offering personally tailored activity sessions to people with dementia in long-term care homes may slightly reduce challenging behavior. Medications No medications have been shown to prevent or cure dementia. Medications may be used to treat the behavioral and cognitive symptoms, but have no effect on the underlying disease process. Acetylcholinesterase inhibitors, such as donepezil, may be useful for Alzheimer's disease, Parkinson's disease dementia, DLB, or vascular dementia. The quality of the evidence is poor and the benefit is small. No difference has been shown between the agents in this family. In a minority of people side effects include a slow heart rate and fainting. Rivastigmine is recommended for treating symptoms in Parkinson's disease dementia. Medications that have anticholinergic effects increase all-cause mortality in people with dementia, although the effect of these medications on cognitive function remains uncertain, according to a systematic review published in 2021. Before prescribing antipsychotic medication in the elderly, an assessment for an underlying cause of the behavior is needed. Severe and life-threatening reactions occur in almost half of people with DLB, and can be fatal after a single dose. People with Lewy body dementias who take neuroleptics are at risk for neuroleptic malignant syndrome, a life-threatening illness. Extreme caution is required in the use of antipsychotic medication in people with DLB because of their sensitivity to these agents. Antipsychotic drugs are used to treat dementia only if non-drug therapies have not worked, and the person's actions threaten themselves or others. Aggressive behavior changes are sometimes the result of other solvable problems, that could make treatment with antipsychotics unnecessary. Because people with dementia can be aggressive, resistant to their treatment, and otherwise disruptive, sometimes antipsychotic drugs are considered as a therapy in response. These drugs have risky adverse effects, including increasing the person's chance of stroke and death. Given these adverse events and small benefit antipsychotics are avoided whenever possible. Generally, stopping antipsychotics for people with dementia does not cause problems, even in those who have been on them a long time. N-methyl-D-aspartate (NMDA) receptor blockers such as memantine may be of benefit but the evidence is less conclusive than for AChEIs. Due to their differing mechanisms of action memantine and acetylcholinesterase inhibitors can be used in combination however the benefit is slight. An extract of Ginkgo biloba known as EGb 761 has been widely used for treating mild to moderate dementia and other neuropsychiatric disorders. Its use is approved throughout Europe. The World Federation of Biological Psychiatry guidelines lists EGb 761 with the same weight of evidence (level B) given to acetylcholinesterase inhibitors, and memantine. EGb 761 is the only one that showed improvement of symptoms in both AD and vascular dementia. EGb 761 is seen as being able to play an important role either on its own or as an add-on particularly when other therapies prove ineffective. EGb 761 is seen to be neuroprotective; it is a free radical scavenger, improves mitochondrial function, and modulates serotonin and dopamine levels. Many studies of its use in mild to moderate dementia have shown it to significantly improve cognitive function, activities of daily living, neuropsychiatric symptoms, and quality of life. However, its use has not been shown to prevent the progression of dementia. While depression is frequently associated with dementia, the use of antidepressants such as selective serotonin reuptake inhibitors (SSRIs) do not appear to affect outcomes. However, the SSRIs sertraline and citalopram have been demonstrated to reduce symptoms of agitation, compared to placebo. No solid evidence indicates that folate or vitamin B12 improves outcomes in those with cognitive problems. Statins have no benefit in dementia. Medications for other health conditions may need to be managed differently for a person who has a dementia diagnosis. It is unclear whether blood pressure medication and dementia are linked. People may experience an increase in cardiovascular-related events if these medications are withdrawn. The Medication Appropriateness Tool for Comorbid Health Conditions in Dementia (MATCH-D) criteria can help identify ways that a diagnosis of dementia changes medication management for other health conditions. These criteria were developed because people with dementia live with an average of five other chronic diseases, which are often managed with medications. The systematic review that informed the criteria were published subsequently in 2018 and updated in 2022. Sleep disturbances Over 40% of people with dementia report sleep problems. Approaches to treating these sleep problems include medications and non-pharmacological approaches. The use of medications to alleviate sleep disturbances that people with dementia often experience has not been well researched, even for medications that are commonly prescribed. In 2012 the American Geriatrics Society recommended that benzodiazepines such as diazepam, and non-benzodiazepine hypnotics, be avoided for people with dementia due to the risks of increased cognitive impairment and falls. Benzodiazepines are also known to promote delirium. Additionally, little evidence supports the effectiveness of benzodiazepines in this population. No clear evidence shows that melatonin or ramelteon improves sleep for people with dementia due to Alzheimer's, but it is used to treat REM sleep behavior disorder in dementia with Lewy bodies. Limited evidence suggests that a low dose of trazodone may improve sleep, however more research is needed. Non-pharmacological approaches have been suggested for treating sleep problems for those with dementia, however, there is no strong evidence or firm conclusions on the effectiveness of different types of interventions, especially for those who are living in an institutionalized setting such as a nursing home or long-term care home. Pain As people age, they experience more health problems, and most health problems associated with aging carry a substantial burden of pain; therefore, between 25% and 50% of older adults experience persistent pain. Seniors with dementia experience the same prevalence of conditions likely to cause pain as seniors without dementia. Pain is often overlooked in older adults and, when screened for, is often poorly assessed, especially among those with dementia, since they become incapable of informing others of their pain. Beyond the issue of humane care, unrelieved pain has functional implications. Persistent pain can lead to decreased ambulation, depressed mood, sleep disturbances, impaired appetite, and exacerbation of cognitive impairment and pain-related interference with activity is a factor contributing to falls in the elderly. Although persistent pain in people with dementia is difficult to communicate, diagnose, and treat, failure to address persistent pain has profound functional, psychosocial and quality of life implications for this vulnerable population. Health professionals often lack the skills and usually lack the time needed to recognize, accurately assess and adequately monitor pain in people with dementia. Family members and friends can make a valuable contribution to the care of a person with dementia by learning to recognize and assess their pain. Educational resources and observational assessment tools are available. Eating difficulties Persons with dementia may have difficulty eating. Whenever it is available as an option, the recommended response to eating problems is having a caretaker assist them. A secondary option for people who cannot swallow effectively is to consider gastrostomy feeding tube placement as a way to give nutrition. However, in bringing comfort and maintaining functional status while lowering risk of aspiration pneumonia and death, assistance with oral feeding is at least as good as tube feeding. Tube-feeding is associated with agitation, increased use of physical and chemical restraints and worsening pressure ulcers. Tube feedings may cause fluid overload, diarrhea, abdominal pain, local complications, less human interaction and may increase the risk of aspiration. Benefits in those with advanced dementia has not been shown. The risks of using tube feeding include agitation, rejection by the person (pulling out the tube, or otherwise physical or chemical immobilization to prevent them from doing this), or developing pressure ulcers. The procedure is directly related to a 1% fatality rate with a 3% major complication rate. The percentage of people at end of life with dementia using feeding tubes in the US has dropped from 12% in 2000 to 6% as of 2014. The immediate and long-term effects of modifying the thickness of fluids for swallowing difficulties in people with dementia are not well known. While thickening fluids may have an immediate positive effect on swallowing and improving oral intake, the long-term impact on the health of the person with dementia should also be considered. Exercise Exercise programs may improve the ability of people with dementia to perform daily activities, but the best type of exercise is still unclear. Getting more exercise can slow the development of cognitive problems such as dementia, proving to reduce the risk of Alzheimer's disease by about 50%. A balance of strength exercise, to help muscles pump blood to the brain, and balance exercises are recommended for aging people. A suggested amount of about hours per week can reduce risks of cognitive decay as well as other health risks like falling. Assistive technology There is a lack of high-quality evidence to determine whether assistive technology effectively supports people with dementia to manage memory issues. Some of the specific things that are used today that helps with dementia today are: clocks, communication aids, electrical appliances the use monitoring, GPS location/ tracking devices, home care robots, in-home cameras, and medication management are just to name a few. Technology has the potential to be a valuable intervention for alleviating loneliness and promoting social connections, supported by available evidence. Alternative medicine Evidence of the therapeutic values of aromatherapy and massage is unclear. It is not clear if cannabinoids are harmful or effective for people with dementia. Palliative care Given the progressive and terminal nature of dementia, palliative care can be helpful to patients and their caregivers by helping people with the disorder and their caregivers understand what to expect, deal with loss of physical and mental abilities, support the person's wishes and goals including surrogate decision making, and discuss wishes for or against CPR and life support. Because the decline can be rapid, and because most people prefer to allow the person with dementia to make their own decisions, palliative care involvement before the late stages of dementia is recommended. Further research is required to determine the appropriate palliative care interventions and how well they help people with advanced dementia. Person-centered care helps maintain the dignity of people with dementia. Remotely delivered information for caregivers Remotely delivered interventions including support, training and information may reduce the burden for the informal caregiver and improve their depressive symptoms. There is no certain evidence that they improve health-related quality of life. In several localities in Japan, digital surveillance may be made available to family members, if a dementia patient is prone to wandering and going missing. Epidemiology The number of cases of dementia worldwide in 2021 was estimated at 55 million, with close to 10 million new cases each year. According to a report by the World Health Organization, "In 2021, Alzheimer’s disease and other forms of dementia ranked as the seventh leading cause of death, killing 1.8 million lives." By 2050, the number of people living with dementia is estimated to be over 150 million globally. Around 7% of people over the age of 65 have dementia, with slightly higher rates (up to 10% of those over 65) in places with relatively high life expectancy. An estimated 58% of people with dementia are living in low and middle income countries. The prevalence of dementia differs in different world regions, ranging from 4.7% in Central Europe to 8.7% in North Africa/Middle East; the prevalence in other regions is estimated to be between 5.6 and 7.6%. The number of people living with dementia is estimated to double every 20 years. In 2016 dementia resulted in about 2.4 million deaths, up from 0.8 million in 1990. The genetic and environmental risk factors for dementia disorders vary by ethnicity. For instance, Alzheimer's disease among Hispanic/Latino and African American subjects exhibit lower risks associated with gene changes in the apolipoprotein E gene than do non-Hispanic white subjects. The annual incidence of dementia diagnosis is nearly 10 million worldwide. Almost half of new dementia cases occur in Asia, followed by Europe (25%), the Americas (18%) and Africa (8%). The incidence of dementia increases exponentially with age, doubling with every 6.3-year increase in age. Dementia affects 5% of the population older than 65 and 20–40% of those older than 85. Rates are slightly higher in women than men at ages 65 and greater. The disease trajectory is varied and the median time from diagnosis to death depends strongly on age at diagnosis, from 6.7 years for people diagnosed aged 60–69 to 1.9 years for people diagnosed at 90 or older. Dementia impacts not only individuals with dementia, but also their carers and the wider society. Among people aged 60 years and over, dementia is ranked the 9th most burdensome condition according to the 2010 Global Burden of Disease (GBD) estimates. The global costs of dementia was around US$818 billion in 2015, a 35.4% increase from US$604 billion in 2010. A new 2024 study reveals that deaths from dementia in the U.S. have tripled in the past 21 years, rising from around 150,000 in 1999 to over 450,000 in 2020; the likelihood of dying from dementia increased across all demographic groups studied. Affected ages About 3% of people between the ages of 65–74 have dementia, 19% between 75 and 84, and nearly half of those over 85 years of age. As more people are living longer, dementia is becoming more common. For people of a specific age, however, it may be becoming less frequent in the developed world, due to a decrease in modifiable risk factors made possible by greater financial and educational resources. It is one of the most common causes of disability among the elderly but can develop before the age of 65 when it is known as early-onset dementia or presenile dementia. Less than 1% of those with Alzheimer's have gene mutations that cause a much earlier development of the disease, around the age of 45, known as early-onset Alzheimer's disease. More than 95% of people with Alzheimer's disease have the sporadic form (late onset, 80–90 years of age). Worldwide the cost of dementia in 2015 was put at US$818 billion. People with dementia are often physically or chemically restrained to a greater degree than necessary, raising issues of human rights. Social stigma is commonly perceived by those with the condition, and also by their caregivers. History Until the end of the 19th century, dementia was a much broader clinical concept. It included mental illness and any type of psychosocial incapacity, including reversible conditions. Dementia at this time simply referred to anyone who had lost the ability to reason, and was applied equally to psychosis, "organic" diseases like syphilis that destroy the brain, and to the dementia associated with old age, which was attributed to "hardening of the arteries". Dementia has been referred to in medical texts since antiquity. One of the earliest known allusions to dementia is attributed to the 7th-century BC Greek philosopher Pythagoras, who divided the human lifespan into six distinct phases: 0–6 (infancy), 7–21 (adolescence), 22–49 (young adulthood), 50–62 (middle age), 63–79 (old age), and 80–death (advanced age). The last two he described as the "senium", a period of mental and physical decay, and that the final phase was when "the scene of mortal existence closes after a great length of time that very fortunately, few of the human species arrive at, where the mind is reduced to the imbecility of the first epoch of infancy". In 550 BC, the Athenian statesman and poet Solon argued that the terms of a man's will might be invalidated if he exhibited loss of judgement due to advanced age. Chinese medical texts made allusions to the condition as well, and the characters for "dementia" translate literally to "foolish old person". Athenian philosophers Aristotle and Plato discussed the mental decline that can come with old age and predicted that this affects everyone who becomes old and nothing can be done to stop this decline from taking place. Plato specifically talked about how the elderly should not be in positions that require responsibility because, "There is not much acumen of the mind that once carried them in their youth, those characteristics one would call judgement, imagination, power of reasoning, and memory. They see them gradually blunted by deterioration and can hardly fulfill their function." For comparison, the Roman statesman Cicero held a view much more in line with modern-day medical wisdom that loss of mental function was not inevitable in the elderly and "affected only those old men who were weak-willed". He spoke of how those who remained mentally active and eager to learn new things could stave off dementia. However, Cicero's views on aging, although progressive, were largely ignored in a world that would be dominated for centuries by Aristotle's medical writings. Physicians during the Roman Empire, such as Galen and Celsus, simply repeated the beliefs of Aristotle while adding few new contributions to medical knowledge. Byzantine physicians sometimes wrote of dementia. It is recorded that at least seven emperors whose lifespans exceeded 70 years displayed signs of cognitive decline. In Constantinople, special hospitals housed those diagnosed with dementia or insanity, but these did not apply to the emperors, who were above the law and whose health conditions could not be publicly acknowledged. Otherwise, little is recorded about dementia in Western medical texts for nearly 1700 years. One of the few references was the 13th-century friar Roger Bacon, who viewed old age as divine punishment for original sin. Although he repeated existing Aristotelian beliefs that dementia was inevitable, he did make the progressive assertion that the brain was the center of memory and thought rather than the heart. Poets, playwrights, and other writers made frequent allusions to the loss of mental function in old age. William Shakespeare notably mentions it in plays such as Hamlet and King Lear. During the 19th century, doctors generally came to believe that elderly dementia was the result of cerebral atherosclerosis, although opinions fluctuated between the idea that it was due to blockage of the major arteries supplying the brain or small strokes within the vessels of the cerebral cortex. In 1907, Bavarian psychiatrist Alois Alzheimer was the first to identify and describe the characteristics of progressive dementia in the brain of 51-year-old Auguste Deter. Deter had begun to behave uncharacteristically, including accusing her husband of adultery, neglecting household chores, exhibiting difficulties writing and engaging in conversations, heightened insomnia, and loss of directional sense. At one point, Deter was reported to have "dragged a bed sheet outside, wandered around wildly, and cried for hours at midnight." Alzheimer began treating Deter when she entered a Frankfurt mental hospital on November 25, 1901. During her ongoing treatment, Deter and her husband struggled to afford the cost of the medical care, and Alzheimer agreed to continue her treatment in exchange for Deter's medical records and donation of her brain upon death. Deter died on April 8, 1906, after succumbing to sepsis and pneumonia. Alzheimer conducted the brain biopsy using the Bielschowsky stain method, which was a new development at the time, and he observed senile plaques, neurofibrillary tangles, and atherosclerotic alteration. At the time, the consensus among medical doctors had been that senile plaques were generally found in older patients, and the occurrence of neurofibrillary tangles was an entirely new observation at the time. Alzheimer presented his findings at the 37th psychiatry conference of southwestern Germany in Tübingen on April 11, 1906; however, the information was poorly received by his peers. By 1910, Alois Alzheimer's teacher, Emil Kraepelin, published a book in which he coined the term "Alzheimer's disease" in an attempt to acknowledge the importance of Alzheimer's discovery. By the 1960s, the link between neurodegenerative diseases and age-related cognitive decline had become more established. By the 1970s, the medical community maintained that vascular dementia was rarer than previously thought and Alzheimer's disease caused the vast majority of old age mental impairments. More recently however, it is believed that dementia is often a mixture of conditions. In 1976, neurologist Robert Katzmann suggested a link between senile dementia and Alzheimer's disease. Katzmann suggested that much of the senile dementia occurring (by definition) after the age of 65, was pathologically identical with Alzheimer's disease occurring in people under age 65 and therefore should not be treated differently. Katzmann thus suggested that Alzheimer's disease, if taken to occur over age 65, is actually common, not rare, and was the fourth- or 5th-leading cause of death, even though rarely reported on death certificates in 1976. A helpful finding was that although the incidence of Alzheimer's disease increased with age (from 5–10% of 75-year-olds to as many as 40–50% of 90-year-olds), no threshold was found by which age all persons developed it. This is shown by documented supercentenarians (people living to 110 or more) who experienced no substantial cognitive impairment. Some evidence suggests that dementia is most likely to develop between ages 80 and 84 and individuals who pass that point without being affected have a lower chance of developing it. Women account for a larger percentage of dementia cases than men. This can be attributed in part to their longer overall lifespan and greater odds of attaining an age where the condition is likely to occur. Much like other diseases associated with aging, dementia was comparatively rare before the 20th century, because few people lived past 80. Conversely, syphilitic dementia was widespread in the developed world until it was largely eradicated by the use of penicillin after World War II. With significant increases in life expectancy thereafter, the number of people over 65 started rapidly climbing. While elderly persons constituted an average of 3–5% of the population prior to 1945, by 2010 many countries reached 10–14% and in Germany and Japan, this figure exceeded 20%. Public awareness of Alzheimer's Disease greatly increased in 1994 when former US president Ronald Reagan announced that he had been diagnosed with the condition. In the 21st century, other types of dementia were differentiated from Alzheimer's disease and vascular dementias (the most common types). This differentiation is on the basis of pathological examination of brain tissues, by symptomatology, and by different patterns of brain metabolic activity in nuclear medical imaging tests such as SPECT and PET scans of the brain. The various forms have differing prognoses and differing epidemiologic risk factors. The main cause for many diseases, including Alzheimer's disease, remains unclear. Terminology Dementia in the elderly was once called senile dementia or senility, and viewed as a normal and somewhat inevitable aspect of aging. By 1913–20 the term dementia praecox was introduced to suggest the development of senile-type dementia at a younger age. Eventually the two terms fused, so that until 1952 physicians used the terms dementia praecox (precocious dementia) and schizophrenia interchangeably. Since then, science has determined that dementia and schizophrenia are two different disorders, though they share some similarities. The term precocious dementia for a mental illness suggested that a type of mental illness like schizophrenia (including paranoia and decreased cognitive capacity) could be expected to arrive normally in all persons with greater age (see paraphrenia). After about 1920, the beginning use of dementia for what is now understood as schizophrenia and senile dementia helped limit the word's meaning to "permanent, irreversible mental deterioration". This began the change to the later use of the term. In recent studies, researchers have seen a connection between those diagnosed with schizophrenia and patients who are diagnosed with dementia, finding a positive correlation between the two diseases. The view that dementia must always be the result of a particular disease process led for a time to the proposed diagnosis of "senile dementia of the Alzheimer's type" (SDAT) in persons over the age of 65, with "Alzheimer's disease" diagnosed in persons younger than 65 who had the same pathology. Eventually, however, it was agreed that the age limit was artificial, and that Alzheimer's disease was the appropriate term for persons with that particular brain pathology, regardless of age. After 1952, mental illnesses including schizophrenia were removed from the category of organic brain syndromes, and thus (by definition) removed from possible causes of "dementing illnesses" (dementias). At the same, however, the traditional cause of senile dementia – "hardening of the arteries" – now returned as a set of dementias of vascular cause (small strokes). These were now termed multi-infarct dementias or vascular dementias. Society and culture The societal cost of dementia is high, especially for caregivers. According to a UK-based study, almost two out of three carers of people with dementia feel lonely. Most of the carers in the study were family members or friends. , the annual cost per Alzheimer's patient in the United States was around $19,144.36. The total costs for the nation is estimated to be about $167.74 billion. By 2030, it is predicted the annual socioeconomic cost will total to about $507 billion, and by 2050 that number is expected to reach $1.89 trillion. This steady increase will be seen not just within the United States but globally. Global estimates for the costs of dementia were $957.56 billion in 2015, but by 2050 the estimated global cost is $9.12 trillion. Many countries consider the care of people living with dementia a national priority and invest in resources and education to better inform health and social service workers, unpaid caregivers, relatives and members of the wider community. Several countries have authored national plans or strategies. These plans recognize that people can live reasonably with dementia for years, as long as the right support and timely access to a diagnosis are available. Former British Prime Minister David Cameron described dementia as a "national crisis", affecting 800,000 people in the United Kingdom. In fact, dementia has become the leading cause of death for women in England. There, as with all mental disorders, people with dementia could potentially be a danger to themselves or others, they can be detained under the Mental Health Act 1983 for assessment, care and treatment. This is a last resort, and is usually avoided by people with family or friends who can ensure care. Some hospitals in Britain work to provide enriched and friendlier care. To make the hospital wards calmer and less overwhelming to residents, staff replaced the usual nurses' station with a collection of smaller desks, similar to a reception area. The incorporation of bright lighting helps increase positive mood and allow residents to see more easily. Driving with dementia can lead to injury or death. Doctors should advise appropriate testing on when to quit driving. The United Kingdom DVLA (Driver & Vehicle Licensing Agency) states that people with dementia who specifically have poor short-term memory, disorientation, or lack of insight or judgment are not allowed to drive, and in these instances the DVLA must be informed so that the driving license can be revoked. They acknowledge that in low-severity cases and those with an early diagnosis, drivers may be permitted to continue driving. Many support networks are available to people with dementia and their families and caregivers. Charitable organizations aim to raise awareness and campaign for the rights of people living with dementia. Support and guidance are available on assessing testamentary capacity in people with dementia. In 2015, Atlantic Philanthropies announced a $177 million gift aimed at understanding and reducing dementia. The recipient was Global Brain Health Institute, a program co-led by the University of California, San Francisco and Trinity College Dublin. This donation is the largest non-capital grant Atlantic has ever made, and the biggest philanthropic donation in Irish history. In October 2020, the Caretaker's last music release, Everywhere at the End of Time, was popularized by TikTok users for its depiction of the stages of dementia. Caregivers were in favor of this phenomenon; Leyland Kirby, the creator of the record, echoed this sentiment, explaining it could cause empathy among a younger public. On November 2, 2020, Scottish billionaire Sir Tom Hunter donated £1 million to dementia charities, after watching a former music teacher with dementia, Paul Harvey, playing one of his own compositions on the piano in a viral video. The donation was announced to be split between the Alzheimer's Society and Music for Dementia. Awareness Celebrities have used their platforms to raise awareness for the different forms of dementia and the need for further support, including former First Lady of California Maria Shriver, Maria Shriver \| My Brain™ \| Alzheimer’s Association Academy Award Winning actor Samuel L. Jackson, Editor-in-Chief of ELLE Magazine Nina Garcia, professional skateboarder Tony Hawk, and others. Additional Alzheimer's awareness has been raised through the diagnoses of high-profile persons themselves, including Actor Bruce Willis Actor Robin Williams Activist Rosa Parks 40th President of the United States, Ronald Reagan Former Mrs. Colorado Springs Joanna Fix Mrs. Colorado Springs uses title, young-onset Alzheimer's diagnosis to spread awareness about dementia TV Host Wendy Williams Wendy Williams diagnosed with same form of dementia as Bruce Willis Musician Tony Bennett Musician Maureen McGovern Dancer and pin-up model Rita Hayworth Rita Hayworth's misdiagnosed struggle	Biology and health sciences	Non-infectious disease	null
38393	https://en.wikipedia.org/wiki/Natural%20rubber	Natural rubber	Rubber, also called India rubber, latex, Amazonian rubber, caucho, or caoutchouc, as initially produced, consists of polymers of the organic compound isoprene, with minor impurities of other organic compounds. Types of polyisoprene that are used as natural rubbers are classified as elastomers. Currently, rubber is harvested mainly in the form of the latex from the Pará rubber tree (Hevea brasiliensis) or others. The latex is a sticky, milky and white colloid drawn off by making incisions in the bark and collecting the fluid in vessels in a process called "tapping". Manufacturers refine this latex into the rubber that is ready for commercial processing. Natural rubber is used extensively in many applications and products, either alone or in combination with other materials. In most of its useful forms, it has a large stretch ratio and high resilience and also is buoyant and water-proof. Industrial demand for rubber-like materials began to outstrip natural rubber supplies by the end of the 19th century, leading to the synthesis of synthetic rubber in 1909 by chemical means. Thailand, Malaysia, Indonesia, and Cambodia are four of the leading rubber producers. Varieties Amazonian rubber tree (Hevea brasiliensis) The major commercial source of natural rubber latex is the Amazonian rubber tree (Hevea brasiliensis), a member of the spurge family, Euphorbiaceae. Once native to Brazil, the species is now pan-tropical. This species is preferred because it grows well under cultivation. A properly managed tree responds to wounding by producing more latex for several years. Congo rubber (Landolphia owariensis and L. spp.) Congo rubber, formerly a major source of rubber, which motivated the atrocities in the Congo Free State, came from vines in the genus Landolphia (L. kirkii, L. heudelotis, and L. owariensis). Dandelion Dandelion milk contains latex. The latex exhibits the same quality as the natural rubber from rubber trees. In the wild types of dandelion, latex content is low and varies greatly. In Nazi Germany, research projects tried to use dandelions as a base for rubber production, but failed. In 2013, by inhibiting one key enzyme and using modern cultivation methods and optimization techniques, scientists in the Fraunhofer Institute for Molecular Biology and Applied Ecology (IME) in Germany developed a cultivar of the Kazakh dandelion (Taraxacum kok-saghyz) that is suitable for commercial production of natural rubber. In collaboration with Continental Tires, IME began a pilot facility. Other Many other plants produce forms of latex rich in isoprene polymers, though not all produce usable forms of polymer as easily as the Pará. Some of them require more elaborate processing to produce anything like usable rubber, and most are more difficult to tap. Some produce other desirable materials, for example gutta-percha (Palaquium gutta) and chicle from Manilkara species. Others that have been commercially exploited, or at least showed promise as rubber sources, include the rubber fig (Ficus elastica), Panama rubber tree (Castilla elastica), various spurges (Euphorbia spp.), lettuce (Lactuca species), the related Scorzonera tau-saghyz, various Taraxacum species, including common dandelion (Taraxacum officinale) and Kazakh dandelion, and, perhaps most importantly for its hypoallergenic properties, guayule (Parthenium argentatum). The term gum rubber is sometimes applied to the tree-obtained version of natural rubber in order to distinguish it from the synthetic version. History The first use of rubber was by the indigenous cultures of Mesoamerica. The earliest archeological evidence of the use of natural latex from the Hevea tree comes from the Olmec culture, in which rubber was first used for making balls for the Mesoamerican ballgame. Rubber was later used by the Maya and Aztec cultures: in addition to making balls, Aztecs used rubber for other purposes, such as making containers and making textiles waterproof by impregnating them with the latex sap. Charles Marie de La Condamine is credited with introducing samples of rubber to the Académie Royale des Sciences of France in 1736. In 1751, he presented a paper by François Fresneau to the Académie (published in 1755) that described many of rubber's properties. This has been referred to as the first scientific paper on rubber. In England, Joseph Priestley, in 1770, observed that a piece of the material was extremely good for rubbing off pencil marks on paper, hence the name "rubber". It slowly made its way around England. In 1764, François Fresnau discovered that turpentine was a rubber solvent. Giovanni Fabbroni is credited with the discovery of naphtha as a rubber solvent in 1779. Charles Goodyear redeveloped vulcanization in 1839, although Mesoamericans had used stabilized rubber for balls and other objects as early as 1600 BC. South America remained the main source of latex rubber used during much of the 19th century. The rubber trade was heavily controlled by business interests but no laws expressly prohibited the export of seeds or plants. In 1876, Henry Wickham smuggled 70,000 Amazonian rubber tree seeds from Brazil and delivered them to Kew Gardens, England. Only 2,400 of these germinated. Seedlings were then sent to India, British Ceylon (Sri Lanka), Dutch East Indies (Indonesia), Singapore, and British Malaya. Malaya (now Peninsular Malaysia) was later to become the biggest producer of rubber. In the early 1900s, the Congo Free State in Africa was also a significant source of natural rubber latex, mostly gathered by forced labor. King Leopold II's colonial state brutally enforced production quotas due to the high price of natural rubber at the time. Tactics to enforce the rubber quotas included removing the hands of victims to prove they had been killed. Soldiers often came back from raids with baskets full of chopped-off hands. Villages that resisted were razed to encourage better compliance locally. The rubber boom in the Amazon also similarly affected indigenous populations to varying degrees. Correrias, or slave raids were frequent in Colombia, Peru and Bolivia where many were either captured or killed. The best-known case of atrocities generated from rubber extraction in South America came from the Putumayo genocide. Between the 1880s and 1913, Julio César Arana and his company, which would become the Peruvian Amazon Company, controlled the Putumayo river. W. E. Hardenburg, Benjamin Saldaña Rocca and Roger Casement were influential figures in exposing these atrocities. Roger Casement was also prominent in revealing the Congo atrocities to the world. Days before entering Iquitos by boat Casement wrote "'Caoutchouc was first called 'india rubber,' because it came from the Indies, and the earliest European use of it was to rub out or erase. It is now called India rubber because it rubs out or erases the Indians." In India, commercial cultivation was introduced by British planters, although the experimental efforts to grow rubber on a commercial scale were initiated as early as 1873 at the Calcutta Botanical Garden. The first commercial Hevea plantations were established at Thattekadu in Kerala in 1902. In later years the plantation expanded to Karnataka, Tamil Nadu and the Andaman and Nicobar Islands of India. Today, India is the world's third-largest producer and fourth-largest consumer of rubber. In Singapore and Malaya, commercial production was heavily promoted by Sir Henry Nicholas Ridley, who served as the first Scientific Director of the Singapore Botanic Gardens from 1888 to 1911. He distributed rubber seeds to many planters and developed the first technique for tapping trees for latex without causing serious harm to the tree. Because of his fervent promotion of this crop, he is popularly remembered by the nickname "Mad Ridley". Pre–World War II Before World War II significant uses included door and window profiles, hoses, belts, gaskets, matting, flooring, and dampeners (antivibration mounts) for the automotive industry. The use of rubber in car tires (initially solid rather than pneumatic) in particular consumed a significant amount of rubber. Gloves (medical, household, and industrial) and toy balloons were large consumers of rubber, although the type of rubber used is concentrated latex. Significant tonnage of rubber was used as adhesives in many manufacturing industries and products, although the two most noticeable were the paper and the carpet industries. Rubber was commonly used to make rubber bands and pencil erasers. Rubber produced as a fiber, sometimes called 'elastic', had significant value to the textile industry because of its excellent elongation and recovery properties. For these purposes, manufactured rubber fiber was made as either an extruded round fiber or rectangular fibers cut into strips from extruded film. Because of its low dye acceptance, feel and appearance, the rubber fiber was either covered by yarn of another fiber or directly woven with other yarns into the fabric. Rubber yarns were used in foundation garments. While rubber is still used in textile manufacturing, its low tenacity limits its use in lightweight garments because latex lacks resistance to oxidizing agents and is damaged by aging, sunlight, oil and perspiration. The textile industry turned to neoprene (polymer of chloroprene), a type of synthetic rubber, as well as another more commonly used elastomer fiber, spandex (also known as elastane), because of their superiority to rubber in both strength and durability. Properties Rubber exhibits unique physical and chemical properties. Rubber's stress–strain behavior exhibits the Mullins effect and the Payne effect and is often modeled as hyperelastic. Rubber strain crystallizes. Because there are weakened allylic C–H bonds in each repeat unit, natural rubber is susceptible to vulcanisation as well as being sensitive to ozone cracking. The two main solvents for rubber are turpentine and naphtha (petroleum). Because rubber does not dissolve easily, the material is finely divided by shredding prior to its immersion. An ammonia solution can be used to prevent the coagulation of raw latex. Rubber begins to melt at approximately . Elasticity On a microscopic scale, relaxed rubber is a disorganized cluster of erratically changing wrinkled chains. In stretched rubber, the chains are almost linear. The restoring force is due to the preponderance of wrinkled conformations over more linear ones. For the quantitative treatment see ideal chain, for more examples see entropic force. Cooling below the glass transition temperature permits local conformational changes but a reordering is practically impossible because of the larger energy barrier for the concerted movement of longer chains. "Frozen" rubber's elasticity is low and strain results from small changes of bond lengths and angles: this caused the Challenger disaster, when the American Space Shuttle's flattened o-rings failed to relax to fill a widening gap. The glass transition is fast and reversible: the force resumes on heating. The parallel chains of stretched rubber are susceptible to crystallization. This takes some time because turns of twisted chains have to move out of the way of the growing crystallites. Crystallization has occurred, for example, when, after days, an inflated toy balloon is found withered at a relatively large remaining volume. Where it is touched, it shrinks because the temperature of the hand is enough to melt the crystals. Vulcanization of rubber creates di- and polysulfide bonds between chains, which limits the degrees of freedom and results in chains that tighten more quickly for a given strain, thereby increasing the elastic force constant and making the rubber harder and less extensible. Malodour Raw rubber storage depots and rubber processing can produce malodour that is serious enough to become a source of complaints and protest to those living in the vicinity. Microbial impurities originate during the processing of block rubber. These impurities break down during storage or thermal degradation and produce volatile organic compounds. Examination of these compounds using gas chromatography/mass spectrometry (GC/MS) and gas chromatography (GC) indicates that they contain sulfur, ammonia, alkenes, ketones, esters, hydrogen sulfide, nitrogen, and low-molecular-weight fatty acids (C2–C5). When latex concentrate is produced from rubber, sulfuric acid is used for coagulation. This produces malodourous hydrogen sulfide. The industry can mitigate these bad odours with scrubber systems. Chemical makeup Rubber is the polymer cis-1,4-polyisoprene – with a molecular weight of 100,000 to 1,000,000 daltons. Typically, a small percentage (up to 5% of dry mass) of other materials, such as proteins, fatty acids, resins, and inorganic materials (salts) are found in natural rubber. Polyisoprene can also be created synthetically, producing what is sometimes referred to as "synthetic natural rubber", but the synthetic and natural routes are distinct. Some natural rubber sources, such as gutta-percha, are composed of trans-1,4-polyisoprene, a structural isomer that has similar properties. Natural rubber is an elastomer and a thermoplastic. Once the rubber is vulcanized, it is a thermoset. Most rubber in everyday use is vulcanized to a point where it shares properties of both; i.e., if it is heated and cooled, it is degraded but not destroyed. The final properties of a rubber item depend not just on the polymer, but also on modifiers and fillers, such as carbon black, factice, whiting and others. Biosynthesis Rubber particles are formed in the cytoplasm of specialized latex-producing cells called laticifers within rubber plants. Rubber particles are surrounded by a single phospholipid membrane with hydrophobic tails pointed inward. The membrane allows biosynthetic proteins to be sequestered at the surface of the growing rubber particle, which allows new monomeric units to be added from outside the biomembrane, but within the lacticifer. The rubber particle is an enzymatically active entity that contains three layers of material, the rubber particle, a biomembrane and free monomeric units. The biomembrane is held tightly to the rubber core by the high negative charge along the double bonds of the rubber polymer backbone. Free monomeric units and conjugated proteins make up the outer layer. The rubber precursor is isopentenyl pyrophosphate (an allylic compound), which elongates by Mg2+-dependent condensation by the action of rubber transferase. The monomer adds to the pyrophosphate end of the growing polymer. The process displaces the terminal high-energy pyrophosphate. The reaction produces a cis polymer. The initiation step is catalyzed by prenyltransferase, which converts three monomers of isopentenyl pyrophosphate into farnesyl pyrophosphate. The farnesyl pyrophosphate can bind to rubber transferase to elongate a new rubber polymer. The required isopentenyl pyrophosphate is obtained from the mevalonate pathway, which derives from acetyl-CoA in the cytosol. In plants, isoprene pyrophosphate can also be obtained from the 1-deox-D-xyulose-5-phosphate/2-C-methyl-D-erythritol-4-phosphate pathway within plasmids. The relative ratio of the farnesyl pyrophosphate initiator unit and isoprenyl pyrophosphate elongation monomer determines the rate of new particle synthesis versus elongation of existing particles. Though rubber is known to be produced by only one enzyme, extracts of latex host numerous small molecular weight proteins with unknown function. The proteins possibly serve as cofactors, as the synthetic rate decreases with complete removal. Production More than 28 million tons of rubber were produced in 2017, of which approximately 47% was natural. Since the bulk is synthetic, which is derived from petroleum, the price of natural rubber is determined, to a large extent, by the prevailing global price of crude oil. Asia was the main source of natural rubber, accounting for about 90% of output in 2021. The three largest producers, Thailand, Indonesia, and Malaysia, together account for around 72% of all natural rubber production. Natural rubber is not cultivated widely in its native continent of South America because of the South American leaf blight, and other natural predators there. Cultivation Rubber latex is extracted from rubber trees. The economic life of rubber trees in plantations is around 32 years, with up to 7 years being an immature phase and about 25 years of productive phase. The soil requirement is well-drained, weathered soil consisting of laterite, lateritic types, sedimentary types, nonlateritic red or alluvial soils. The climatic conditions for optimum growth of rubber trees are: Rainfall of around evenly distributed without any marked dry season and with at least 100 rainy days per year Temperature range of about , with a monthly mean of Atmospheric humidity of around 80% About 2,000 hours sunshine per year at the rate of six hours per day throughout the year Absence of strong winds Many high-yielding clones have been developed for commercial planting. These clones yield more than of dry rubber per year, under ideal conditions. Rubber production has been linked to deforestation. Rubber therefore is one of seven commodities included in the 2023 EU Regulation on Deforestation-free products (EUDR), which aims to guarantee that the products European Union (EU) citizens consume do not contribute to deforestation or forest degradation worldwide. Collection In places such as Kerala and Sri Lanka, where coconuts are in abundance, the half shell of coconut was used as the latex collection container. Glazed pottery or aluminium or plastic cups became more common in Kerala-India and other countries. The cups are supported by a wire that encircles the tree. This wire incorporates a spring so it can stretch as the tree grows. The latex is led into the cup by a galvanised "spout" knocked into the bark. Rubber tapping normally takes place early in the morning, when the internal pressure of the tree is highest. A good tapper can tap a tree every 20 seconds on a standard half-spiral system, and a common daily "task" size is between 450 and 650 trees. Trees are usually tapped on alternate or third days, although many variations in timing, length and number of cuts are used. "Tappers would make a slash in the bark with a small hatchet. These slanting cuts allowed latex to flow from ducts located on the exterior or the inner layer of bark (cambium) of the tree. Since the cambium controls the growth of the tree, growth stops if it is cut. Thus, rubber tapping demanded accuracy, so that the incisions would not be too many given the size of the tree, or too deep, which could stunt its growth or kill it." It is usual to tap a panel at least twice, sometimes three times, during the tree's life. The economic life of the tree depends on how well the tapping is carried out, as the critical factor is bark consumption. A standard in Malaysia for alternate daily tapping is 25 cm (vertical) bark consumption per year. The latex-containing tubes in the bark ascend in a spiral to the right. For this reason, tapping cuts usually ascend to the left to cut more tubes. The trees drip latex for about four hours, stopping as latex coagulates naturally on the tapping cut, thus blocking the latex tubes in the bark. Tappers usually rest and have a meal after finishing their tapping work and then start collecting the liquid "field latex" at about midday. Field coagula The four types of field coagula are "cuplump", "treelace", "smallholders' lump", and "earth scrap". Each has significantly different properties. Some trees continue to drip after the collection leading to a small amount of "cup lump" that is collected at the next tapping. The latex that coagulates on the cut is also collected as "tree lace". Tree lace and cup lump together account for 10%–20% of the dry rubber produced. Latex that drips onto the ground, "earth scrap", is also collected periodically for processing of low-grade product. Cup lump Cup lump is the coagulated material found in the collection cup when the tapper next visits the tree to tap it again. It arises from latex clinging to the walls of the cup after the latex was last poured into the bucket, and from late-dripping latex exuded before the latex-carrying vessels of the tree become blocked. It is of higher purity and of greater value than the other three types. 'Cup lumps' can also be used to describe a completely different type of coagulate that has collected in smallholder plantations over a period of 1–2 weeks. After tapping all of the trees, the tapper will return to each tree and stir in some type of acid, which allows the newly harvested latex to mix with the previously coagulated material. The rubber/acid mixture is what gives rubber plantations, markets, and factories a strong odor. Tree lace Tree lace is the coagulum strip that the tapper peels off the previous cut before making a new cut. It usually has higher copper and manganese contents than cup lump. Both copper and manganese are pro-oxidants and can damage the physical properties of the dry rubber. Smallholders' lump Smallholders' lump is produced by smallholders, who collect rubber from trees far from the nearest factory. Many Indonesian smallholders, who farm paddies in remote areas, tap dispersed trees on their way to work in the paddy fields and collect the latex (or the coagulated latex) on their way home. As it is often impossible to preserve the latex sufficiently to get it to a factory that processes latex in time for it to be used to make high quality products, and as the latex would anyway have coagulated by the time it reached the factory, the smallholder will coagulate it by any means available, in any container available. Some smallholders use small containers, buckets etc., but often the latex is coagulated in holes in the ground, which are usually lined with plastic sheeting. Acidic materials and fermented fruit juices are used to coagulate the latex – a form of assisted biological coagulation. Little care is taken to exclude twigs, leaves, and even bark from the lumps that are formed, which may also include tree lace. Earth scrap Earth scrap is material that gathers around the base of the tree. It arises from latex overflowing from the cut and running down the bark, from rain flooding a collection cup containing latex, and from spillage from tappers' buckets during collection. It contains soil and other contaminants, and has variable rubber content, depending on the amount of contaminants. Earth scrap is collected by field workers two or three times a year and may be cleaned in a scrap-washer to recover the rubber, or sold to a contractor who cleans it and recovers the rubber. It is of low quality. Processing Latex coagulates in the cups if kept for long and must be collected before this happens. The collected latex, "field latex", is transferred into coagulation tanks for the preparation of dry rubber or transferred into air-tight containers with sieving for ammoniation. Ammoniation, invented by patent lawyer and vice-president of the United States Rubber Company Ernest Hopkinson around 1920, preserves the latex in a colloidal state for longer periods of time. Latex is generally processed into either latex concentrate for manufacture of dipped goods or coagulated under controlled, clean conditions using formic acid. The coagulated latex can then be processed into the higher-grade, technically specified block rubbers such as SVR 3L or SVR CV or used to produce Ribbed Smoke Sheet grades. Naturally coagulated rubber (cup lump) is used in the manufacture of TSR10 and TSR20 grade rubbers. Processing for these grades is a size reduction and cleaning process to remove contamination and prepare the material for the final stage of drying. The dried material is then baled and palletized for storage and shipment. Molecular structure Rubber is a natural polymer of isoprene (polyisoprene), and an elastomer (a stretchy polymer). Polymers are simply chains of molecules that can be linked together. Rubber is one of the few naturally occurring polymers and prized for its high stretch ratio, resilience, and water-proof properties. Other examples of natural polymers include tortoise shell, amber, and animal horn. When harvested, latex rubber takes the form of latex, an opaque, white, milky suspension of rubber particles in water. It is then transformed through industrial processes to the solid form widely seen in manufactured goods. Vulcanized rubber Natural rubber is reactive and vulnerable to oxidization, but it can be stabilized through a heating process called vulcanization. Vulcanization is a process by which the rubber is heated and sulfur, peroxide, or bisphenol are added to improve resistance and elasticity and to prevent it from oxidizing. Carbon black, which can be derived from a petroleum refinery or other natural incineration processes, is sometimes used as an additive to rubber to improve its strength, especially in vehicle tires. During vulcanization, rubber's polyisoprene molecules (long chains of isoprene) are heated and cross-linked with molecular bonds to sulfur, forming a 3-D matrix. The optimal percentage of sulfur is approximately 10%. In this form, the polyisoprene molecules orientation is still random but they become aligned when the rubber is stretched. This sulfur vulcanization makes the rubber stronger and more rigid, but still very elastic. And through the vulcanization process, the sulfur and latex are meant to be totally used up in individual form. Transportation Natural rubber latex is shipped from factories in Southeast Asia, South America, and West and Central Africa to destinations around the world. As the cost of natural rubber has risen significantly and rubber products are dense, the shipping methods offering the lowest cost per unit weight are preferred. Depending on destination, warehouse availability, and transportation conditions, some methods are preferred by certain buyers. In international trade, latex rubber is mostly shipped in 20-foot ocean containers. Inside the container, smaller containers are used to store the latex. Rubber shortage and global economics There is growing concern for the future supply of rubber due to various factors, including plant disease, climate change, and the volatile market price of rubber. Producers of natural rubber are mostly small family-held plantations, often serving large industrial aggregators. High volatility in the price of rubber affects rubber plantation investment, and farmers may remove their rubber trees if the international market spot price of a seemingly more profitable crop (for example palm oil) surges in relation to rubber. For instance, during the 2020 and 2021 international COVID-19 pandemic, demand for rubber gloves surged, leading to a spike in rubber prices of about 30%. In addition to the pandemic, demand exceeded supply in part because long term plantations had been torn out and replaced with other crops over the previous 5–10 years, and other areas were affected by climate-fueled natural disasters. In this environment, producers did increase their prices in keeping with supply and demand dynamics, putting upward price pressure on the whole downstream supply chain. Uses Uncured rubber is used for cements; for adhesive, insulating, and friction tapes; and for crepe rubber used in insulating blankets and footwear. Vulcanized rubber has many more applications. Resistance to abrasion makes softer kinds of rubber valuable for the treads of vehicle tires and conveyor belts, and makes hard rubber valuable for pump housings and piping used in the handling of abrasive sludge. The flexibility of rubber is appealing in hoses, tires and rollers for devices ranging from domestic clothes wringers to printing presses; its elasticity makes it suitable for various kinds of shock absorbers and for specialized machinery mountings designed to reduce vibration. Its relative gas impermeability makes it useful in the manufacture of articles such as air hoses, balloons, balls and cushions. The resistance of rubber to water and to the action of most fluid chemicals has led to its use in rainwear, diving gear, and chemical and medicinal tubing and as a lining for storage tanks, processing equipment and railroad tank cars. Because of their electrical resistance, soft rubber goods are used as insulation and for protective gloves, shoes, and blankets; hard rubber is used for articles such as telephone housings and parts for radio sets, meters, and other electrical instruments. The coefficient of friction of rubber, which is high on dry surfaces and low on wet surfaces, leads to its use for power-transmission belting, highly flexible couplings, and for water-lubricated bearings in deep-well pumps. Indian rubber balls or lacrosse balls are made of rubber. Around 25 million tonnes of rubber are produced each year, of which 30 percent is natural. The remainder is synthetic rubber derived from petrochemical sources. The top end of latex production results in latex products such as surgeons' gloves, balloons, and other relatively high-value products. The mid-range which comes from the technically specified natural rubber materials ends up largely in tires but also in conveyor belts, marine products, windshield wipers, and miscellaneous goods. Natural rubber offers good elasticity, while synthetic materials tend to offer better resistance to environmental factors such as oils, temperature, chemicals, and ultraviolet light. "Cured rubber" is rubber that has been compounded and subjected to the vulcanisation process to create cross-links within the rubber matrix. Rubber can be added to cement to improve its properties. Allergic reactions Some people have a serious latex allergy, and exposure to natural latex rubber products such as latex gloves can cause anaphylactic shock. The antigenic proteins found in Hevea latex are reduced by about 99.9 percent (though not eliminated) through vulcanization processing. Latex from non-Hevea sources, such as guayule, can be used without allergic reaction by persons with an allergy to Hevea latex. Some allergic reactions are not to the latex itself, but from residues of chemicals used to accelerate the cross-linking process. Although this may be confused with an allergy to latex, it is distinct from it, typically taking the form of Type IV hypersensitivity in the presence of traces of specific processing chemicals. Microbial degradation Natural rubber is susceptible to degradation by a wide range of bacteria. The bacteria Streptomyces coelicolor, Pseudomonas citronellolis, and Nocardia spp. are capable of degrading vulcanized natural rubber.	Technology	Materials	null
38413	https://en.wikipedia.org/wiki/Activation%20energy	Activation energy	In the Arrhenius model of reaction rates, activation energy is the minimum amount of energy that must be available to reactants for a chemical reaction to occur. The activation energy (Ea) of a reaction is measured in kilojoules per mole (kJ/mol) or kilocalories per mole (kcal/mol). Activation energy can be thought of as the magnitude of the potential barrier (sometimes called the energy barrier) separating minima of the potential energy surface pertaining to the initial and final thermodynamic state. For a chemical reaction to proceed at a reasonable rate, the temperature of the system should be high enough such that there exists an appreciable number of molecules with translational energy equal to or greater than the activation energy. The term "activation energy" was introduced in 1889 by the Swedish scientist Svante Arrhenius. Other uses Although less commonly used, activation energy also applies to nuclear reactions and various other physical phenomena. Temperature dependence and the relation to the Arrhenius equation The Arrhenius equation gives the quantitative basis of the relationship between the activation energy and the rate at which a reaction proceeds. From the equation, the activation energy can be found through the relation where A is the pre-exponential factor for the reaction, R is the universal gas constant, T is the absolute temperature (usually in kelvins), and k is the reaction rate coefficient. Even without knowing A, Ea can be evaluated from the variation in reaction rate coefficients as a function of temperature (within the validity of the Arrhenius equation). At a more advanced level, the net Arrhenius activation energy term from the Arrhenius equation is best regarded as an experimentally determined parameter that indicates the sensitivity of the reaction rate to temperature. There are two objections to associating this activation energy with the threshold barrier for an elementary reaction. First, it is often unclear as to whether or not reaction does proceed in one step; threshold barriers that are averaged out over all elementary steps have little theoretical value. Second, even if the reaction being studied is elementary, a spectrum of individual collisions contributes to rate constants obtained from bulk ('bulb') experiments involving billions of molecules, with many different reactant collision geometries and angles, different translational and (possibly) vibrational energies—all of which may lead to different microscopic reaction rates. Catalysts A substance that modifies the transition state to lower the activation energy is termed a catalyst; a catalyst composed only of protein and (if applicable) small molecule cofactors is termed an enzyme. A catalyst increases the rate of reaction without being consumed in the reaction. In addition, the catalyst lowers the activation energy, but it does not change the energies of the original reactants or products, and so does not change equilibrium. Rather, the reactant energy and the product energy remain the same and only the activation energy is altered (lowered). A catalyst is able to reduce the activation energy by forming a transition state in a more favorable manner. Catalysts, by nature, create a more "comfortable" fit for the substrate of a reaction to progress to a transition state. This is possible due to a release of energy that occurs when the substrate binds to the active site of a catalyst. This energy is known as Binding Energy. Upon binding to a catalyst, substrates partake in numerous stabilizing forces while within the active site (e.g. hydrogen bonding or van der Waals forces). Specific and favorable bonding occurs within the active site until the substrate forms to become the high-energy transition state. Forming the transition state is more favorable with the catalyst because the favorable stabilizing interactions within the active site release energy. A chemical reaction is able to manufacture a high-energy transition state molecule more readily when there is a stabilizing fit within the active site of a catalyst. The binding energy of a reaction is this energy released when favorable interactions between substrate and catalyst occur. The binding energy released assists in achieving the unstable transition state. Reactions without catalysts need a higher input of energy to achieve the transition state. Non-catalyzed reactions do not have free energy available from active site stabilizing interactions, such as catalytic enzyme reactions. Relationship with Gibbs energy of activation In the Arrhenius equation, the term activation energy (Ea) is used to describe the energy required to reach the transition state, and the exponential relationship holds. In transition state theory, a more sophisticated model of the relationship between reaction rates and the transition state, a superficially similar mathematical relationship, the Eyring equation, is used to describe the rate constant of a reaction: . However, instead of modeling the temperature dependence of reaction rate phenomenologically, the Eyring equation models individual elementary steps of a reaction. Thus, for a multistep process, there is no straightforward relationship between the two models. Nevertheless, the functional forms of the Arrhenius and Eyring equations are similar, and for a one-step process, simple and chemically meaningful correspondences can be drawn between Arrhenius and Eyring parameters. Instead of also using Ea, the Eyring equation uses the concept of Gibbs energy and the symbol ΔG‡ to denote the Gibbs energy of activation to achieve the transition state. In the equation, kB and h are the Boltzmann and Planck constants, respectively. Although the equations look similar, it is important to note that the Gibbs energy contains an entropic term in addition to the enthalpic one. In the Arrhenius equation, this entropic term is accounted for by the pre-exponential factor A. More specifically, we can write the Gibbs free energy of activation in terms of enthalpy and entropy of activation: . Then, for a unimolecular, one-step reaction, the approximate relationships and hold. Note, however, that in Arrhenius theory proper, A is temperature independent, while here, there is a linear dependence on T. For a one-step unimolecular process whose half-life at room temperature is about 2 hours, ΔG‡ is approximately 23 kcal/mol. This is also the roughly the magnitude of Ea for a reaction that proceeds over several hours at room temperature. Due to the relatively small magnitude of TΔS‡ and RT at ordinary temperatures for most reactions, in sloppy discourse, Ea, ΔG‡, and ΔH‡ are often conflated and all referred to as the "activation energy". The enthalpy, entropy and Gibbs energy of activation are more correctly written as Δ‡Ho, Δ‡So and Δ‡Go respectively, where the o indicates a quantity evaluated between standard states. However, some authors omit the o in order to simplify the notation. The total free energy change of a reaction is independent of the activation energy however. Physical and chemical reactions can be either exergonic or endergonic, but the activation energy is not related to the spontaneity of a reaction. The overall reaction energy change is not altered by the activation energy. Negative activation energy In some cases, rates of reaction decrease with increasing temperature. When following an approximately exponential relationship so the rate constant can still be fit to an Arrhenius expression, this results in a negative value of Ea. Elementary reactions exhibiting negative activation energies are typically barrierless reactions, in which the reaction proceeding relies on the capture of the molecules in a potential well. Increasing the temperature leads to a reduced probability of the colliding molecules capturing one another (with more glancing collisions not leading to reaction as the higher momentum carries the colliding particles out of the potential well), expressed as a reaction cross section that decreases with increasing temperature. Such a situation no longer leads itself to direct interpretations as the height of a potential barrier. Some multistep reactions can also have apparent negative activation energies. For example, the overall rate constant k for a two-step reaction A B, B → C is given by k = k2K1, where k2 is the rate constant of the rate-limiting slow second step and K1 is the equilibrium constant of the rapid first step. In some reactions, K1 decreases with temperature more rapidly than k2 increases, so that k actually decreases with temperature corresponding to a negative observed activation energy. An example is the oxidation of nitric oxide which is a termolecular reaction 2 NO + O2 -> 2 NO2. The rate law is with a negative activation energy. This is explained by the two-step mechanism: 2 NO <=> N2O2 and N2O2 + O2 -> 2 NO2. Certain cationic polymerization reactions have negative activation energies so that the rate decreases with temperature. For chain-growth polymerization, the overall activation energy is , where i, p and t refer respectively to initiation, propagation and termination steps. The propagation step normally has a very small activation energy, so that the overall value is negative if the activation energy for termination is larger than that for initiation. The normal range of overall activation energies for cationic polymerization varies from .	Physical sciences	Kinetics	Chemistry
38428	https://en.wikipedia.org/wiki/Deer	Deer	A deer (: deer) or true deer is a hoofed ruminant ungulate of the family Cervidae (informally the deer family). Cervidae is divided into subfamilies Cervinae (which includes, among others, muntjac, elk (wapiti), red deer, and fallow deer) and Capreolinae (which includes, among others reindeer (caribou), white-tailed deer, roe deer, and moose). Male deer of almost all species (except the water deer), as well as female reindeer, grow and shed new antlers each year. These antlers are bony extensions of the skull and are often used for combat between males. The musk deer (Moschidae) of Asia and chevrotains (Tragulidae) of tropical African and Asian forests are separate families that are also in the ruminant clade Ruminantia; they are not especially closely related to Cervidae. Deer appear in art from Paleolithic cave paintings onwards, and they have played a role in mythology, religion, and literature throughout history, as well as in heraldry, such as red deer that appear in the coat of arms of Åland. Their economic importance includes the use of their meat as venison, their skins as soft, strong buckskin, and their antlers as handles for knives. Deer hunting has been a popular activity since the Middle Ages and remains a resource for many families today. Etymology and terminology The word deer was originally broad in meaning, becoming more specific with time. Old English and Middle English meant a wild animal of any kind. Cognates of Old English in other dead Germanic languages have the general sense of animal, such as Old High German tior, Old Norse or , Gothic dius, Old Saxon dier, and Old Frisian diar. This general sense gave way to the modern English sense by the end of the Middle English period, around 1500. All modern Germanic languages save English and Scots retain the more general sense: for example, Dutch/Frisian , German , and Norwegian dyr mean . For many types of deer in modern English usage, the male is a buck and the female a doe, but the terms vary with dialect, and according to the size of the species. The male red deer is a stag, while for other large species the male is a bull, the female a cow, as in cattle. In older usage, the male of any species is a hart, especially if over five years old, and the female is a hind, especially if three or more years old. The young of small species is a fawn and of large species a calf; a very small young may be a kid. A castrated male is a havier. A group of any species is a herd. The adjective of relation is cervine; like the family name Cervidae, this is from , meaning or . Distribution Deer live in a variety of biomes, ranging from tundra to the tropical rainforest. While often associated with forests, many deer are ecotone species that live in transitional areas between forests and thickets (for cover) and prairie and savanna (open space). The majority of large deer species inhabit temperate mixed deciduous forest, mountain mixed coniferous forest, tropical seasonal/dry forest, and savanna habitats around the world. Clearing open areas within forests to some extent may actually benefit deer populations by exposing the understory and allowing the types of grasses, weeds, and herbs to grow that deer like to eat. Access to adjacent croplands may also benefit deer. Adequate forest or brush cover must still be provided for populations to grow and thrive. Deer are widely distributed, with indigenous representatives in all continents except Antarctica and Australia, though Africa has only one native deer, the Barbary stag, a subspecies of red deer that is confined to the Atlas Mountains in the northwest of the continent. Another extinct species of deer, Megaceroides algericus, was present in North Africa until 6000 years ago. Fallow deer have been introduced to South Africa. Small species of brocket deer and pudús of Central and South America, and muntjacs of Asia generally occupy dense forests and are less often seen in open spaces, with the possible exception of the Indian muntjac. There are also several species of deer that are highly specialized and live almost exclusively in mountains, grasslands, swamps, and "wet" savannas, or riparian corridors surrounded by deserts. Some deer have a circumpolar distribution in both North America and Eurasia. Examples include the caribou that live in Arctic tundra and taiga (boreal forests) and moose that inhabit taiga and adjacent areas. Huemul deer (taruca and Chilean huemul) of South America's Andes fill the ecological niches of the ibex and wild goat, with the fawns behaving more like goat kids. The highest concentration of large deer species in temperate North America lies in the Canadian Rocky Mountain and Columbia Mountain regions between Alberta and British Columbia where all five North American deer species (white-tailed deer, mule deer, caribou, elk, and moose) can be found. This region has several clusters of national parks including Mount Revelstoke National Park, Glacier National Park (Canada), Yoho National Park, and Kootenay National Park on the British Columbia side, and Banff National Park, Jasper National Park, and Glacier National Park (U.S.) on the Alberta and Montana sides. Mountain slope habitats vary from moist coniferous/mixed forested habitats to dry subalpine/pine forests with alpine meadows higher up. The foothills and river valleys between the mountain ranges provide a mosaic of cropland and deciduous parklands. The rare woodland caribou have the most restricted range living at higher altitudes in the subalpine meadows and alpine tundra areas of some of the mountain ranges. Elk and mule deer both migrate between the alpine meadows and lower coniferous forests and tend to be most common in this region. Elk also inhabit river valley bottomlands, which they share with White-tailed deer. The White-tailed deer have recently expanded their range within the foothills and river valley bottoms of the Canadian Rockies owing to conversion of land to cropland and the clearing of coniferous forests allowing more deciduous vegetation to grow up the mountain slopes. They also live in the aspen parklands north of Calgary and Edmonton, where they share habitat with the moose. The adjacent Great Plains grassland habitats are left to herds of elk, American bison, and pronghorn. The Eurasian Continent (including the Indian Subcontinent) boasts the most species of deer in the world, with most species being found in Asia. Europe, in comparison, has lower diversity in plant and animal species. Many national parks and protected reserves in Europe have populations of red deer, roe deer, and fallow deer. These species have long been associated with the continent of Europe, but also inhabit Asia Minor, the Caucasus Mountains, and Northwestern Iran. "European" fallow deer historically lived over much of Europe during the Ice Ages, but afterwards became restricted primarily to the Anatolian Peninsula, in present-day Turkey. Present-day fallow deer populations in Europe are a result of historic man-made introductions of this species, first to the Mediterranean regions of Europe, then eventually to the rest of Europe. They were initially park animals that later escaped and reestablished themselves in the wild. Historically, Europe's deer species shared their deciduous forest habitat with other herbivores, such as the extinct tarpan (forest horse), extinct aurochs (forest ox), and the endangered wisent (European bison). Good places to see deer in Europe include the Scottish Highlands, the Austrian Alps, the wetlands between Austria, Hungary, and the Czech Republic, and some National Parks, including Doñana National Park in Spain, the Veluwe in the Netherlands, the Ardennes in Belgium, and Białowieża National Park in Poland. Spain, Eastern Europe, and the Caucasus Mountains have forest areas that are not only home to sizable deer populations but also other animals that were once abundant such as the wisent, Eurasian lynx, Iberian lynx, wolves, and brown bears. The highest concentration of large deer species in temperate Asia occurs in the mixed deciduous forests, mountain coniferous forests, and taiga bordering North Korea, Manchuria (Northeastern China), and the Ussuri Region (Russia). These are among some of the richest deciduous and coniferous forests in the world where one can find Siberian roe deer, sika deer, elk, and moose. Asian caribou occupy the northern fringes of this region along the Sino-Russian border. Deer such as the sika deer, Thorold's deer, Central Asian red deer, and elk have historically been farmed for their antlers by Han Chinese, Turkic peoples, Tungusic peoples, Mongolians, and Koreans. Like the Sami people of Finland and Scandinavia, the Tungusic peoples, Mongolians, and Turkic peoples of Southern Siberia, Northern Mongolia, and the Ussuri Region have also taken to raising semi-domesticated herds of Asian caribou. The highest concentration of large deer species in the tropics occurs in Southern Asia in India's Indo-Gangetic Plain Region and Nepal's Terai Region. These fertile plains consist of tropical seasonal moist deciduous, dry deciduous forests, and both dry and wet savannas that are home to chital, hog deer, barasingha, Indian sambar, and Indian muntjac. Grazing species such as the endangered barasingha and very common chital are gregarious and live in large herds. Indian sambar can be gregarious but are usually solitary or live in smaller herds. Hog deer are solitary and have lower densities than Indian muntjac. Deer can be seen in several national parks in India, Nepal, and Sri Lanka of which Kanha National Park, Dudhwa National Park, and Chitwan National Park are most famous. Sri Lanka's Wilpattu National Park and Yala National Park have large herds of Indian sambar and chital. The Indian sambar are more gregarious in Sri Lanka than other parts of their range and tend to form larger herds than elsewhere. The Chao Praya River Valley of Thailand was once primarily tropical seasonal moist deciduous forest and wet savanna that hosted populations of hog deer, the now-extinct Schomburgk's deer, Eld's deer, Indian sambar, and Indian muntjac. Both the hog deer and Eld's deer are rare, whereas Indian sambar and Indian muntjac thrive in protected national parks, such as Khao Yai. Many of these South Asian and Southeast Asian deer species also share their habitat with other herbivores, such as Asian elephants, the various Asian rhinoceros species, various antelope species (such as nilgai, four-horned antelope, blackbuck, and Indian gazelle in India), and wild oxen (such as wild Asian water buffalo, gaur, banteng, and kouprey). One way that different herbivores can survive together in a given area is for each species to have different food preferences, although there may be some overlap. As a result of acclimatisation society releases in the 19th century, Australia has six introduced species of deer that have established sustainable wild populations. They are fallow deer, red deer, sambar, hog deer, rusa, and chital. Red deer were introduced into New Zealand in 1851 from English and Scottish stock. Many have been domesticated in deer farms since the late 1960s and are common farm animals there now. Seven other species of deer were introduced into New Zealand but none are as widespread as red deer. Description Deer constitute the second most diverse family of artiodactyla after bovids. Though of a similar build, deer are strongly distinguished from antelopes by their antlers, which are temporary and regularly regrown unlike the permanent horns of bovids. Characteristics typical of deer include long, powerful legs, a diminutive tail and long ears. Deer exhibit a broad variation in physical proportions. The largest extant deer is the moose, which is nearly tall and weighs up to . The elk stands at the shoulder and weighs . The northern pudu is the smallest deer in the world; it reaches merely at the shoulder and weighs . The southern pudu is only slightly taller and heavier. Sexual dimorphism is quite pronounced – in most species males tend to be larger than females, and, except for the reindeer, only males have antlers. Coat colour generally varies between red and brown, though it can be as dark as chocolate brown in the tufted deer or have a grayish tinge as in elk. Different species of brocket deer vary from gray to reddish brown in coat colour. Several species such as the chital, the fallow deer and the sika deer feature white spots on a brown coat. Coat of reindeer shows notable geographical variation. Deer undergo two moults in a year; for instance, in red deer the red, thin-haired summer coat is gradually replaced by the dense, greyish brown winter coat in autumn, which in turn gives way to the summer coat in the following spring. Moulting is affected by the photoperiod. Deer are also excellent jumpers and swimmers. Deer are ruminants, or cud-chewers, and have a four-chambered stomach. Some deer, such as those on the island of Rùm, do consume meat when it is available. Nearly all deer have a facial gland in front of each eye. The gland contains a strongly scented pheromone, used to mark its home range. Bucks of a wide range of species open these glands wide when angry or excited. All deer have a liver without a gallbladder. Deer also have a tapetum lucidum, which gives them sufficiently good night vision. Antlers All male deer have antlers, with the exception of the water deer, in which males have long tusk-like canines that reach below the lower jaw. Females generally lack antlers, though female reindeer bear antlers smaller and less branched than those of the males. Occasionally females in other species may develop antlers, especially in telemetacarpal deer such as European roe deer, red deer, white-tailed deer and mule deer and less often in plesiometacarpal deer. A study of antlered female white-tailed deer noted that antlers tend to be small and malformed, and are shed frequently around the time of parturition. The fallow deer and the various subspecies of the reindeer have the largest as well as the heaviest antlers, both in absolute terms as well as in proportion to body mass (an average of eight grams per kilogram of body mass); the tufted deer, on the other hand, has the smallest antlers of all deer, while the pudú has the lightest antlers with respect to body mass (0.6 g per kilogram of body mass). The structure of antlers show considerable variation; while fallow deer and elk antlers are palmate (with a broad central portion), white-tailed deer antlers include a series of tines sprouting upward from a forward-curving main beam, and those of the pudú are mere spikes. Antler development begins from the pedicel, a bony structure that appears on the top of the skull by the time the animal is a year old. The pedicel gives rise to a spiky antler the following year, that is replaced by a branched antler in the third year. This process of losing a set of antlers to develop a larger and more branched set continues for the rest of the life. The antlers emerge as soft tissues (known as velvet antlers) and progressively harden into bony structures (known as hard antlers), following mineralisation and blockage of blood vessels in the tissue, from the tip to the base. Antlers might be one of the most exaggerated male secondary sexual characteristics, and are intended primarily for reproductive success through sexual selection and for combat. The tines (forks) on the antlers create grooves that allow another male's antlers to lock into place. This allows the males to wrestle without risking injury to the face. Antlers are correlated to an individual's position in the social hierarchy and its behaviour. For instance, the heavier the antlers, the higher the individual's status in the social hierarchy, and the greater the delay in shedding the antlers; males with larger antlers tend to be more aggressive and dominant over others. Antlers can be an honest signal of genetic quality; males with larger antlers relative to body size tend to have increased resistance to pathogens and higher reproductive capacity. In elk in Yellowstone National Park, antlers also provide protection against predation by wolves. Homology of tines, that is, the branching structure of antlers among species, have been discussed before the 1900s. Recently, a new method to describe the branching structure of antlers and determining homology of tines was developed. Teeth Most deer bear 32 teeth; the corresponding dental formula is: . The elk and the reindeer may be exceptions, as they may retain their upper canines and thus have 34 teeth (dental formula: ). The Chinese water deer, tufted deer, and muntjac have enlarged upper canine teeth forming sharp tusks, while other species often lack upper canines altogether. The cheek teeth of deer have crescent ridges of enamel, which enable them to grind a wide variety of vegetation. The teeth of deer are adapted to feeding on vegetation, and like other ruminants, they lack upper incisors, instead having a tough pad at the front of their upper jaw. Biology Diet Deer are browsers, and feed primarily on foliage of grasses, sedges, forbs, shrubs and trees, secondarily on lichens in northern latitudes during winter. They have small, unspecialized stomachs by ruminant standards, and high nutrition requirements. Rather than eating and digesting vast quantities of low-grade fibrous food as, for example, sheep and cattle do, deer select easily digestible shoots, young leaves, fresh grasses, soft twigs, fruit, fungi, and lichens. The low-fibered food, after minimal fermentation and shredding, passes rapidly through the alimentary canal. The deer require a large amount of minerals such as calcium and phosphate in order to support antler growth, and this further necessitates a nutrient-rich diet. There are some reports of deer engaging in carnivorous activity, such as eating dead alewives along lakeshores or depredating the nests of northern bobwhites. Reproduction Nearly all cervids are so-called uniparental species: the young, known in most species as fawns, are only cared for by the mother, most often called a doe. A doe generally has one or two fawns at a time (triplets, while not unknown, are uncommon). Mating season typically begins in later August and lasts until December. Some species mate until early March. The gestation period is anywhere up to ten months for the European roe deer. Most fawns are born with their fur covered with white spots, though in many species they lose these spots by the end of their first winter. In the first twenty minutes of a fawn's life, the fawn begins to take its first steps. Its mother licks it clean until it is almost free of scent, so predators will not find it. Its mother leaves often to graze, and the fawn does not like to be left behind. Sometimes its mother must gently push it down with her foot. The fawn stays hidden in the grass for one week until it is strong enough to walk with its mother. The fawn and its mother stay together for about one year. A male usually leaves and never sees his mother again, but females sometimes come back with their own fawns and form small herds. Disease In some areas of the UK, deer (especially fallow deer due to their gregarious behaviour) have been implicated as a possible reservoir for transmission of bovine tuberculosis, a disease which in the UK in 2005 cost £90 million in attempts to eradicate. In New Zealand, deer are thought to be important as vectors picking up M. bovis in areas where brushtail possums Trichosurus vulpecula are infected, and transferring it to previously uninfected possums when their carcasses are scavenged elsewhere. The white-tailed deer Odocoileus virginianus has been confirmed as the sole maintenance host in the Michigan outbreak of bovine tuberculosis which remains a significant barrier to the US nationwide eradication of the disease in livestock. Moose and deer can carry rabies. Docile moose may suffer from brain worm, a helminth which drills holes through the brain in its search for a suitable place to lay its eggs. A government biologist states that "They move around looking for the right spot and never really find it." Deer appear to be immune to this parasite; it passes through the digestive system and is excreted in the feces. The parasite is not screened by the moose intestine, and passes into the brain where damage is done that is externally apparent, both in behaviour and in gait. Deer, elk and moose in North America may suffer from chronic wasting disease, which was identified at a Colorado laboratory in the 1960s and is believed to be a prion disease. Out of an abundance of caution hunters are advised to avoid contact with specified risk material (SRM) such as the brain, spinal column or lymph nodes. Deboning the meat when butchering and sanitizing the knives and other tools used to butcher are amongst other government recommendations. Evolution Deer are believed to have evolved from antlerless, tusked ancestors that resembled modern duikers and diminutive deer in the early Eocene, and gradually developed into the first antlered cervoids (the superfamily of cervids and related extinct families) in the Miocene. Eventually, with the development of antlers, the tusks as well as the upper incisors disappeared. Thus, evolution of deer took nearly 30 million years. Biologist Valerius Geist suggests evolution to have occurred in stages. There are not many prominent fossils to trace this evolution, but only fragments of skeletons and antlers that might be easily confused with false antlers of non-cervid species. Eocene The ruminants, ancestors of the Cervidae, are believed to have evolved from Diacodexis, the earliest known artiodactyl (even-toed ungulate), 50–55 Mya in the Eocene. Diacodexis, nearly the size of a rabbit, featured the talus bone characteristic of all modern even-toed ungulates. This ancestor and its relatives occurred throughout North America and Eurasia, but were on the decline by at least 46 Mya. Analysis of a nearly complete skeleton of Diacodexis discovered in 1982 gave rise to speculation that this ancestor could be closer to the non-ruminants than the ruminants. Andromeryx is another prominent prehistoric ruminant, but appears to be closer to the tragulids. Oligocene The formation of the Himalayas and the Alps brought about significant geographic changes. This was the chief reason behind the extensive diversification of deer-like forms and the emergence of cervids from the Oligocene to the early Pliocene. The latter half of the Oligocene (28–34 Mya) saw the appearance of the European Eumeryx and the North American Leptomeryx. The latter resembled modern-day bovids and cervids in dental morphology (for instance, it had brachyodont molars), while the former was more advanced. Other deer-like forms included the North American Blastomeryx and the European Dremotherium; these sabre-toothed animals are believed to have been the direct ancestors of all modern antlered deer, though they themselves lacked antlers. Another contemporaneous form was the four-horned protoceratid Protoceras, that was replaced by Syndyoceras in the Miocene; these animals were unique in having a horn on the nose. Late Eocene fossils dated approximately 35 million years ago, which were found in North America, show that Syndyoceras had bony skull outgrowths that resembled non-deciduous antlers. Miocene Fossil evidence suggests that the earliest members of the superfamily Cervoidea appeared in Eurasia in the Miocene. Dicrocerus, Euprox and Heteroprox were probably the first antlered cervids. Dicrocerus featured single-forked antlers that were shed regularly. Stephanocemas had more developed and diffuse ("crowned") antlers. Procervulus (Palaeomerycidae) also had antlers that were not shed. Contemporary forms such as the merycodontines eventually gave rise to the modern pronghorn. The Cervinae emerged as the first group of extant cervids around 7–9 Mya, during the late Miocene in central Asia. The tribe Muntiacini made its appearance as † Muntiacus leilaoensis around 7–8 Mya; The early muntjacs varied in size–as small as hares or as large as fallow deer. They had tusks for fighting and antlers for defence. Capreolinae followed soon after; Alceini appeared 6.4–8.4 Mya. Around this period, the Tethys Ocean disappeared to give way to vast stretches of grassland; these provided the deer with abundant protein-rich vegetation that led to the development of ornamental antlers and allowed populations to flourish and colonise areas. As antlers had become pronounced, the canines were either lost or became poorly represented (as in elk), probably because diet was no longer browse-dominated and antlers were better display organs. In muntjac and tufted deer, the antlers as well as the canines are small. The tragulids have long canines to this day. Pliocene With the onset of the Pliocene, the global climate became cooler. A fall in the sea-level led to massive glaciation; consequently, grasslands abounded in nutritious forage. Thus a new spurt in deer populations ensued. The oldest member of Cervini, † Cervocerus novorossiae, appeared around the transition from Miocene to Pliocene (4.2–6 Mya) in Eurasia; cervine fossils from early Pliocene to as late as the Pleistocene have been excavated in China and the Himalayas. While Cervus and Dama appeared nearly 3 Mya, Axis emerged during the late Pliocene–Pleistocene. The tribes Capreolini and Rangiferini appeared around 4–7 Mya. Around 5 Mya, the rangiferina † Bretzia and † Eocoileus were the first cervids to reach North America. This implies the Bering Strait could be crossed during the late Miocene–Pliocene; this appears highly probable as the camelids migrated into Asia from North America around the same time. Deer invaded South America in the late Pliocene (2.5–3 Mya) as part of the Great American Interchange, thanks to the recently formed Isthmus of Panama, and emerged successful due to the small number of competing ruminants in the continent. Pleistocene Large deer with impressive antlers evolved during the early Pleistocene, probably as a result of abundant resources to drive evolution. The early Pleistocene cervid † Eucladoceros was comparable in size to the modern elk. † Megaloceros (Pliocene–Pleistocene) featured the Irish elk (M. giganteus), one of the largest known cervids. The Irish elk reached at the shoulder and had heavy antlers that spanned from tip to tip. These large animals were traditionally thought to have faced extinction due to conflict between sexual selection for large antlers and body and natural selection for a smaller form, but a combination of anthropogenic and climatic pressures is now thought to be the most likely culprit. Meanwhile, the moose and reindeer radiated into North America from Siberia. Taxonomy and classification Deer constitute the artiodactyl family Cervidae. This family was first described by German zoologist Georg August Goldfuss in Handbuch der Zoologie (1820). Three subfamilies were recognised: Capreolinae (first described by the English zoologist Joshua Brookes in 1828), Cervinae (described by Goldfuss) and Hydropotinae (first described by French zoologist Édouard Louis Trouessart in 1898). Other attempts at the classification of deer have been based on morphological and genetic differences. The Anglo-Irish naturalist Victor Brooke suggested in 1878 that deer could be bifurcated into two classes on the according to the features of the second and fifth metacarpal bones of their forelimbs: Plesiometacarpalia (most Old World deer) and Telemetacarpalia (most New World deer). He treated the musk deer as a cervid, placing it under Telemetacarpalia. While the telemetacarpal deer showed only those elements located far from the joint, the plesiometacarpal deer retained the elements closer to the joint as well. Differentiation on the basis of diploid number of chromosomes in the late 20th century has been flawed by several inconsistencies. In 1987, the zoologists Colin Groves and Peter Grubb identified three subfamilies: Cervinae, Hydropotinae and Odocoileinae; they noted that the hydropotines lack antlers, and the other two subfamilies differ in their skeletal morphology. They reverted from this classification in 2000. Molecular phylogenetic analyses since the latter half of the 2000s all show that hydropotes is a sister taxon of Capreolus, and “Hydropotinae” became outdated subfamily. External relationships Until 2003, it was understood that the family Moschidae (musk deer) was sister to Cervidae. Then a phylogenetic study by Alexandre Hassanin (of National Museum of Natural History, France) and colleagues, based on mitochondrial and nuclear analyses, revealed that Moschidae and Bovidae form a clade sister to Cervidae. According to the study, Cervidae diverged from the Bovidae-Moschidae clade 27 to 28 million years ago. The following cladogram is based on the 2003 study. Internal relationships A 2006 phylogenetic study of the internal relationships in Cervidae by Clément Gilbert and colleagues divided the family into two major clades: Capreolinae (telemetacarpal or New World deer) and Cervinae (plesiometacarpal or Old World deer). Studies in the late 20th century suggested a similar bifurcation in the family. This as well as previous studies support monophyly in Cervinae, while Capreolinae appears paraphyletic. The 2006 study identified two lineages in Cervinae, Cervini (comprising the genera Axis, Cervus, Dama and Rucervus) and Muntiacini (Muntiacus and Elaphodus). Capreolinae featured three lineages, Alceini (Alces species), Capreolini (Capreolus and the subfamily Hydropotinae) and Rangiferini (Blastocerus, Hippocamelus, Mazama, Odocoileus, Pudu and Rangifer species). The following cladogram is based on the 2006 study. Human interaction Prehistoric Deer were an important source of food for early hominids. In China, Homo erectus fed upon the sika deer, while the red deer was hunted in Germany. In the Upper Palaeolithic, the reindeer was the staple food for Cro-Magnon people, while the cave paintings at Lascaux in southwestern France include some 90 images of stags. In China, deer continued to be a main source of food for millennia even after people began farming, and it is possible that sika and other deer benefited from the frequently abandoned field sites. Historic Deer had a central role in the ancient art, culture and mythology of various peoples including the Hittites, the ancient Egyptians, the Celts, the ancient Greeks, and certain East Asian cultures. For instance, the Stag Hunt Mosaic of ancient Pella, under the Kingdom of Macedonia (4th century BC), possibly depicts Alexander the Great hunting a deer with Hephaestion. In Japanese Shintoism, the sika deer is believed to be a messenger to the gods. In China, deer are associated with great medicinal significance; deer penis is thought by some in China to have aphrodisiac properties. Spotted deer are believed in China to accompany the god of longevity. Deer was the principal sacrificial animal for the Huichal Indians of Mexico. In medieval Europe, deer appeared in hunting scenes and coats-of-arms. Deer are depicted in many materials by various pre-Hispanic civilizations in the Andes. The common male given name Oscar is taken from the Irish Language, where it is derived from two elements: the first, os, means "deer"; the second element, cara, means "friend". The name is borne by a famous hero of Irish mythology—Oscar, grandson of Fionn Mac Cumhail. The name was popularised in the 18th century by James Macpherson, creator of 'Ossianic poetry'. Literary Deer have been an integral part of fables and other literary works since the inception of writing. Stags were used as symbols in the latter Sumerian writings. For instance, the boat of Sumerian god Enki is named the Stag of Azbu. There are several mentions of the animal in the Rigveda as well as the Bible. In the Indian epic Ramayana, Sita is lured by a golden deer which Rama tries to catch. In the absence of both Rama and Lakshman, Ravana kidnaps Sita. Many of the allegorical Aesop's fables, such as "The Stag at the Pool", "The One-Eyed Doe" and "The Stag and a Lion", personify deer to give moral lessons. For instance, "The Sick Stag" gives the message that uncaring friends can do more harm than good. The Yaqui deer song accompanies the deer dance which is performed by a pascola [from the Spanish 'pascua', Easter] dancer (also known as a deer dancer). Pascolas would perform at religious and social functions many times of the year, especially during Lent and Easter. In one of Rudolf Erich Raspe's 1785 stories of Baron Munchausen's Narrative of his Marvellous Travels and Campaigns in Russia, the baron encounters a stag while eating cherries and, without ammunition, fires the cherry-pits at the stag with his musket, but it escapes. The next year, the baron encounters a stag with a cherry tree growing from its head; presumably this is the animal he had shot at the previous year. In Christmas lore (such as in the narrative poem "A Visit from St. Nicholas"), reindeer are often depicted pulling the sleigh of Santa Claus. Marjorie Kinnan Rawlings's Pulitzer Prize-winning 1938 novel The Yearling was about a boy's relationship with a baby deer. The fiction book Fire Bringer is about a young fawn who goes on a quest to save the Herla, the deer kind. In the 1942 Walt Disney Pictures film, Bambi is a white-tailed deer, while in Felix Salten's original 1923 book Bambi, a Life in the Woods, he is a roe deer. In C. S. Lewis's 1950 fantasy novel The Lion, the Witch and the Wardrobe the adult Pevensies, now kings and queens of Narnia, chase the White Stag on a hunt, as the Stag is said to grant its captor a wish. The hunt is key in returning the Pevensies to their home in England. In the 1979 book The Animals of Farthing Wood, The Great White Stag is the leader of all the animals. Heraldic Deer of various types appear frequently in European heraldry. In the British armory, the term "stag" is typically used to refer to antlered male red deer, while "buck" indicates an antlered male fallow deer. Stags and bucks appear in a number of attitudes, referred to as "lodged" when the deer is lying down, "trippant" when it has one leg raised, "courant" when it is running, "springing" when in the act of leaping, "statant" when it is standing with all hooves on the ground and looking ahead, and "at gaze" when otherwise statant but looking at the viewer. Stags' heads are also frequently used; these are typically portrayed without an attached neck and as facing the viewer, in which case they are termed "caboshed". Examples of deer in coats of arms can be found in the arms of Hertfordshire, England, and its county town of Hertford; both are examples of canting arms. A deer appears on the arms of the Israeli Postal Authority. Coats of arms featuring deer include those of Baden-Württemberg, Dotternhausen, Thierachern, Friolzheim, Bauen, Albstadt, and Dassel in Germany; of the Earls Bathurst in England; of Balakhna, Gusev, Nizhny Novgorod, Odintsovo, Slavsk and Yamalo-Nenets in Russia; of Berezhany, Sambir in Ukraine; of Åland, Finland; of Gjemnes, Hitra, Hjartdal and Rendalen in Norway; of Jelenia Góra, Poland; of Umeå, Sweden; of Queensland, Australia; of Cervera, Catalonia; of Selonia and Semigallia in Latvia; and of Chile. Other types of deer used in heraldry include the hind, portrayed much like the stag or buck but without antlers, as well as the reindeer and winged stags. Winged stags are used as supporters in the arms of the de Carteret family. The sea-stag, having the antlers, head, forelegs and upper body of a stag and the tail of a mermaid, is often found in German heraldry. Economic Deer have long had economic significance to humans. Deer meat, known as venison, is highly nutritious. Due to the inherently wild nature and diet of deer, venison is most often obtained through deer hunting. In the United States, it is produced in small amounts compared to beef, but still represents a significant trade. Deer hunting is a popular activity in the U.S. that can provide the hunter's family with high quality meat and generates revenue for states and the federal government from the sales of licenses, permits and tags. The 2006 survey by the U.S. Fish and Wildlife Service estimates that license sales generate approximately $700 million annually. This revenue generally goes to support conservation efforts in the states where the licenses are purchased. Overall, the U.S. Fish and Wildlife Service estimates that big game hunting for deer and elk generates approximately $11.8 billion annually in hunting-related travel, equipment and related expenditures. Conservation laws prevent the sale of unlicensed wild game meat, although it may be donated. Deer have often been bred in captivity as ornaments for parks, but only in the case of reindeer has thorough domestication succeeded. By 2012, some 25,000 tons of red deer were raised on farms in North America. The Sami of Scandinavia and the Kola Peninsula of Russia and other nomadic peoples of northern Asia use reindeer for food, clothing, and transport. Others are bred for hunting are selected based on the size of the antlers. The major deer-producing countries are New Zealand, the market leader, with Ireland, Great Britain and Germany. The trade earns over $100 million annually for these countries. Automobile collisions with deer can impose a significant cost on the economy. In the U.S., about 1.5 million deer-vehicle collisions occur each year, according to the National Highway Traffic Safety Administration. Those accidents cause about 150 human deaths and $1.1 billion in property damage annually. In Scotland, several roads including the A82, the A87 and the A835 have had significant enough problems with deer vehicle collisions (DVCs) that sets of vehicle activated automatic warning signs have been installed along these roads. The skins make a peculiarly strong, soft leather, known as buckskin. There is nothing special about skins with the fur still on since the hair is brittle and soon falls off. The hooves and antlers are used for ornamental purposes, especially the antlers of the roe deer, which are utilized for making umbrella handles, and for similar purposes; elk antlers is often employed in making knife handles. Among the Inuit, the traditional ulu women's knife was made with an antler or ivory handle. In China, a traditional chinese medicine is made from stag antler, and the antlers of certain species are eaten when "in the velvet". Antlers can also be boiled down to release the protein gelatin, which is used as a topical treatment for skin irritation and is also used in cooking. Since the early 20th century, deer have become commonly thought of as pests in New Zealand due to a lack of predators on the island causing population numbers to increase and begin encroaching on more populated areas. They compete with livestock for resources, as well as cause excess erosion and wreak havoc on wild plant species and agriculture alike. They can also have an effect on the conservation efforts of other plant and animal species, as they can critically offset the balance within an environment by drastically depleting diversity within forests.	Biology and health sciences	Artiodactyla	null
38431	https://en.wikipedia.org/wiki/Wool	Wool	Wool is the textile fiber obtained from sheep and other mammals, especially goats, rabbits, and camelids. The term may also refer to inorganic materials, such as mineral wool and glass wool, that have some properties similar to animal wool. As an animal fiber, wool consists of protein together with a small percentage of lipids. This makes it chemically quite distinct from cotton and other plant fibers, which are mainly cellulose. Characteristics Wool is produced by follicles which are small cells located in the skin. These follicles are located in the upper layer of the skin called the epidermis and push down into the second skin layer called the dermis as the wool fibers grow. Follicles can be classed as either primary or secondary follicles. Primary follicles produce three types of fiber: kemp, medullated fibers, and true wool fibers. Secondary follicles only produce true wool fibers. Medullated fibers share nearly identical characteristics to hair and are long but lack crimp and elasticity. Kemp fibers are very coarse and shed out. Wool's crimp refers to the strong natural wave present in each wool fiber as it is presented on the animal. Wool's crimp, and to a lesser degree scales, make it easier to spin the fleece by helping the individual fibers attach, so they stay together. Because of the crimp, wool fabrics have greater bulk than other textiles, and they hold air, which causes the fabric to retain heat. Wool has a high specific thermal resistance, so it impedes heat transfer in general. This effect has benefited desert peoples, as Bedouins and Tuaregs use wool clothes for insulation. The felting of wool occurs upon hammering or other mechanical agitation as the microscopic barbs on the surface of wool fibers hook together. Felting generally comes under two main areas, dry felting and wet felting. Wet felting occurs when water and a lubricant (especially an alkali such as soap) are applied to the wool which is then agitated until the fibers mix and bond together. Temperature shock while damp or wet accentuates the felting process. Some natural felting can occur on the animal's back. Wool has several qualities that distinguish it from hair or fur: it is crimped and elastic. The amount of crimp corresponds to the fineness of the wool fibers. A fine wool like Merino may have up to 40 crimps per centimetre (100 crimps per inch), while coarser wool like karakul may have less than one (one or two crimps per inch). In contrast, hair has little if any scale and no crimp, and little ability to bind into yarn. On sheep, the hair part of the fleece is called kemp. The relative amounts of kemp to wool vary from breed to breed and make some fleeces more desirable for spinning, felting, or carding into batts for quilts or other insulating products, including the famous tweed cloth of Scotland. Wool fibers readily absorb moisture, but are not hollow. Wool can absorb almost one-third of its own weight in water. Wool absorbs sound like many other fabrics. It is generally a creamy white color, although some breeds of sheep produce natural colors, such as black, brown, silver, and random mixes. Wool ignites at a higher temperature than cotton and some synthetic fibers. It has a lower rate of flame spread, a lower rate of heat release, a lower heat of combustion, and does not melt or drip; it forms a char that is insulating and self-extinguishing, and it contributes less to toxic gases and smoke than other flooring products when used in carpets. Wool carpets are specified for high safety environments, such as trains and aircraft. Wool is usually specified for garments for firefighters, soldiers, and others in occupations where they are exposed to the likelihood of fire. Wool causes an allergic reaction in some people. Processing Shearing Sheep shearing is the process in which a worker (a shearer) cuts off the woollen fleece of a sheep. After shearing, wool-classers separate the wool into four main categories: fleece (which makes up the vast bulk) broken bellies locks The quality of fleeces is determined by a technique known as wool classing, whereby a qualified person, called a wool classer, groups wools of similar grading together to maximize the return for the farmer or sheep owner. In Australia, before being auctioned, all Merino fleece wool is objectively measured for average diameter (micron), yield (including the amount of vegetable matter), staple length, staple strength, and sometimes color and comfort factor. Scouring Wool straight off a sheep is known as "raw wool", "greasy wool" or "wool in the grease". This wool contains a high level of valuable lanolin, as well as the sheep's dead skin and sweat residue, and generally also contains pesticides and vegetable matter from the animal's environment. Before the wool can be used for commercial purposes, it must be scoured, a process of cleaning the greasy wool. Scouring may be as simple as a bath in warm water or as complicated as an industrial process using detergent and alkali in specialized equipment. In north west England, special potash pits were constructed to produce potash used in the manufacture of a soft soap for scouring locally produced white wool. Vegetable matter in commercial wool is often removed by chemical carbonization. In less-processed wools, vegetable matter may be removed by hand and some of the lanolin left intact through the use of gentler detergents. This semigrease wool can be worked into yarn and knitted into particularly water-resistant mittens or sweaters, such as those of the Aran Island fishermen. Lanolin removed from wool is widely used in cosmetic products such as hand creams. Fineness and yield Raw wool has many impurities; vegetable matter, sand, dirt and yolk which is a mixture of suint (sweat), grease, urine stains and dung locks. The sheep's body yields many types of wool with differing strengths, thicknesses, length of staple and impurities. The raw wool (greasy) is processed into 'top'. 'Worsted top' requires strong straight and parallel fibres. The quality of wool is determined by its fiber diameter, crimp, yield, color, and staple strength. Fiber diameter is the single most important wool characteristic determining quality and price. Merino wool is typically in length and is very fine (between 12 and 24 microns). The finest and most valuable wool comes from Merino hoggets. Wool taken from sheep produced for meat is typically coarser, and has fibers in length. Damage or breaks in the wool can occur if the sheep is stressed while it is growing its fleece, resulting in a thin spot where the fleece is likely to break. Wool is also separated into grades based on the measurement of the wool's diameter in microns and also its style. These grades may vary depending on the breed or purpose of the wool. For example: Any wool finer than 25 microns can be used for garments, while coarser grades are used for outerwear or rugs. The finer the wool, the softer it is, while coarser grades are more durable and less prone to pilling. The finest Australian and New Zealand Merino wools are known as 1PP, which is the industry benchmark of excellence for Merino wool 16.9 microns and finer. This style represents the top level of fineness, character, color, and style as determined on the basis of a series of parameters in accordance with the original dictates of British wool as applied by the Australian Wool Exchange (AWEX) Council. Only a few dozen of the millions of bales auctioned every year can be classified and marked 1PP. In the United States, three classifications of wool are named in the Wool Products Labeling Act of 1939. Wool is "the fiber from the fleece of the sheep or lamb or hair of the Angora or Cashmere goat (and may include the so-called specialty fibers from the hair of the camel, alpaca, llama, and vicuna) which has never been reclaimed from any woven or felted wool product". "Virgin wool" and "new wool" are also used to refer to such never used wool. There are two categories of recycled wool (also called reclaimed or shoddy wool). "Reprocessed wool" identifies "wool which has been woven or felted into a wool product and subsequently reduced to a fibrous state without having been used by the ultimate consumer". "Reused wool" refers to such wool that has been used by the ultimate consumer. History Wild sheep were more hairy than woolly. Although sheep were domesticated some 9,000 to 11,000 years ago, archaeological evidence from statuary found at sites in Iran suggests selection for woolly sheep may have begun around 6000 BC, with the earliest woven wool garments having only been dated to two to three thousand years later. Woolly sheep were introduced into Europe from the Near East in the early part of the 4th millennium BC. The oldest known European wool textile, , was preserved in a Danish bog. Prior to invention of shears—probably in the Iron Age—the wool was plucked out by hand or by bronze combs. In Roman times, wool, linen, and leather clothed the European population; cotton from India was a curiosity of which only naturalists had heard, and silks, imported along the Silk Road from China, were extravagant luxury goods. Pliny the Elder records in his Natural History that the reputation for producing the finest wool was enjoyed by Tarentum, where selective breeding had produced sheep with superior fleeces, but which required special care. In medieval times, as trade connections expanded, the Champagne fairs revolved around the production of wool cloth in small centers such as Provins. The network developed by the annual fairs meant the woolens of Provins might find their way to Naples, Sicily, Cyprus, Majorca, Spain, and even Constantinople. The wool trade developed into serious business, a generator of capital. In the 13th century, the wool trade became the economic engine of the Low Countries and central Italy. By the end of the 14th century, Italy predominated. The Florentine wool guild, Arte della Lana, sent the imported English wool to the San Martino convent for processing. Italian wool from Abruzzo and Spanish merino wools were processed at Garbo workshops. Abruzzo wool had once been the most accessible for the Florentine guild, until improved relations with merchants in Iberia made merino wool more available. By the 16th century Italian wool exports to the Levant had declined, eventually replaced by silk production. The value of exports of English raw wool were rivaled only by the 15th-century sheepwalks of Castile and were a significant source of income to the English crown, which in 1275 had imposed an export tax on wool called the "Great Custom". The importance of wool to the English economy can be seen in the fact that since the 14th century, the presiding officer of the House of Lords has sat on the "Woolsack", a chair stuffed with wool. Economies of scale were instituted in the Cistercian houses, which had accumulated great tracts of land during the 12th and early 13th centuries, when land prices were low and labor still scarce. Raw wool was baled and shipped from North Sea ports to the textile cities of Flanders, notably Ypres and Ghent, where it was dyed and worked up as cloth. At the time of the Black Death, English textile industries consumed about 10% of English wool production. The English textile trade grew during the 15th century, to the point where export of wool was discouraged. Over the centuries, various British laws controlled the wool trade or required the use of wool even in burials. The smuggling of wool out of the country, known as owling, was at one time punishable by the cutting off of a hand. After the Restoration, fine English woolens began to compete with silks in the international market, partly aided by the Navigation Acts; in 1699, the English crown forbade its American colonies to trade wool with anyone but England herself. A great deal of the value of woollen textiles was in the dyeing and finishing of the woven product. In each of the centers of the textile trade, the manufacturing process came to be subdivided into a collection of trades, overseen by an entrepreneur in a system called by the English the "putting-out" system, or "cottage industry", and the Verlagssystem by the Germans. In this system of producing wool cloth, once perpetuated in the production of Harris tweeds, the entrepreneur provides the raw materials and an advance, the remainder being paid upon delivery of the product. Written contracts bound the artisans to specified terms. Fernand Braudel traces the appearance of the system in the 13th-century economic boom, quoting a document of 1275. The system effectively bypassed the guilds' restrictions. Before the flowering of the Renaissance, the Medici and other great banking houses of Florence had built their wealth and banking system on their textile industry based on wool, overseen by the Arte della Lana, the wool guild: wool textile interests guided Florentine policies. Francesco Datini, the "merchant of Prato", established in 1383 an Arte della Lana for that small Tuscan city. The sheepwalks of Castile were controlled by the Mesta union of sheep owners. They shaped the landscape and the fortunes of the meseta that lies in the heart of the Iberian peninsula; in the 16th century, a unified Spain allowed export of Merino lambs only with royal permission. The German wool market – based on sheep of Spanish origin – did not overtake British wool until comparatively late. Later, the Industrial Revolution introduced mass production technology into wool and wool cloth manufacturing. Australia's colonial economy was based on sheep raising, and the Australian wool trade eventually overtook that of the Germans by 1845, furnishing wool for Bradford, which developed as the heart of industrialized woolens production. Due to decreasing demand with increased use of synthetic fibers, wool production is much less than what it was in the past. The collapse in the price of wool began in late 1966 with a 40% drop; with occasional interruptions, the price has tended down. The result has been sharply reduced production and movement of resources into production of other commodities, in the case of sheep growers, to production of meat. Superwash wool (or washable wool) technology first appeared in the early 1970s to produce wool that has been specially treated so it is machine washable and may be tumble-dried. This wool is produced using an acid bath that removes the "scales" from the fiber, or by coating the fiber with a polymer that prevents the scales from attaching to each other and causing shrinkage. This process results in a fiber that holds longevity and durability over synthetic materials, while retaining its shape. In December 2004, a bale of the then world's finest wool, averaging 11.8 microns, sold for AU$3,000 per kilogram at auction in Melbourne. This fleece wool tested with an average yield of 74.5%, long, and had 40 newtons per kilotex strength. The result was A$279,000 for the bale. The finest bale of wool ever auctioned was sold for a seasonal record of AU$2690 per kilo during June 2008. This bale was produced by the Hillcreston Pinehill Partnership and measured 11.6 microns, 72.1% yield, and had a 43 newtons per kilotex strength measurement. The bale realized $247,480 and was exported to India. In 2007, a new wool suit was developed and sold in Japan that can be washed in the shower, and which dries off ready to wear within hours with no ironing required. The suit was developed using Australian Merino wool, and it enables woven products made from wool, such as suits, trousers, and skirts, to be cleaned using a domestic shower at home. In December 2006, the General Assembly of the United Nations proclaimed 2009 to be the International Year of Natural Fibres, so as to raise the profile of wool and other natural fibers. Production Global wool production is about per year, of which 60% goes into apparel. Wool comprises ca 3% of the global textile market, but its value is higher owing to dyeing and other modifications of the material. Australia is a leading producer of wool which is mostly from Merino sheep but has been eclipsed by China in terms of total weight. New Zealand (2016) is the third-largest producer of wool, and the largest producer of crossbred wool. Breeds such as Lincoln, Romney, Drysdale, and Elliotdale produce coarser fibers, and wool from these sheep is usually used for making carpets. In the United States, Texas, New Mexico, and Colorado have large commercial sheep flocks and their mainstay is the Rambouillet (or French Merino). Also, a thriving home-flock contingent of small-scale farmers raise small hobby flocks of specialty sheep for the hand-spinning market. These small-scale farmers offer a wide selection of fleece. Global woolclip (total amount of wool shorn) 2020 Organic wool has gained in popularity. This wool is limited in supply and much of it comes from New Zealand and Australia. Organic wool has become easier to find in clothing and other products, but these products often carry a higher price. Wool is environmentally preferable (as compared to petroleum-based nylon or polypropylene) as a material for carpets, as well, in particular when combined with a natural binding and the use of formaldehyde-free glues. Animal rights groups have noted issues with the production of wool, such as mulesing. Marketing Australia About 85% of wool sold in Australia is sold by open cry auction. Other countries The British Wool Marketing Board operates a central marketing system for UK fleece wool with the aim of achieving the best possible net returns for farmers. Less than half of New Zealand's wool is sold at auction, while around 45% of farmers sell wool directly to private buyers and end-users. United States sheep producers market wool with private or cooperative wool warehouses, but wool pools are common in many states. In some cases, wool is pooled in a local market area, but sold through a wool warehouse. Wool offered with objective measurement test results is preferred. Imported apparel wool and carpet wool goes directly to central markets, where it is handled by the large merchants and manufacturers. Yarn Shoddy or recycled wool is made by cutting or tearing apart existing wool fabric and respinning the resulting fibers. As this process makes the wool fibers shorter, the remanufactured fabric is inferior to the original. The recycled wool may be mixed with raw wool, wool noil, or another fiber such as cotton to increase the average fiber length. Such yarns are typically used as weft yarns with a cotton warp. This process was invented in the Heavy Woollen District of West Yorkshire and created a microeconomy in this area for many years. Worsted is a strong, long-staple, combed wool yarn with a hard surface. Woolen is a soft, short-staple, carded wool yarn typically used for knitting. In traditional weaving, woolen weft yarn (for softness and warmth) is frequently combined with a worsted warp yarn for strength on the loom. Uses In addition to clothing, wool has been used for blankets, suits, horse rugs, saddle cloths, carpeting, insulation and upholstery. Dyed wool can be used to create other forms of art such as wet and needle felting. Wool felt covers piano hammers, and it is used to absorb odors and noise in heavy machinery and stereo speakers. Ancient Greeks lined their helmets with felt, and Roman legionnaires used breastplates made of wool felt. Wool as well as cotton has also been traditionally used for cloth diapers. Wool fiber exteriors are hydrophobic (repel water) and the interior of the wool fiber is hygroscopic (attracts water); this makes a wool garment suitable cover for a wet diaper by inhibiting wicking, so outer garments remain dry. Wool felted and treated with lanolin is water resistant, air permeable, and slightly antibacterial, so it resists the buildup of odor. Some modern cloth diapers use felted wool fabric for covers, and there are several modern commercial knitting patterns for wool diaper covers. Initial studies of woollen underwear have found it prevented heat and sweat rashes because it more readily absorbs the moisture than other fibers. As an animal protein, wool can be used as a soil fertilizer, being a slow-release source of nitrogen. Researchers at the Royal Melbourne Institute of Technology school of fashion and textiles have discovered a blend of wool and Kevlar, the synthetic fiber widely used in body armor, was lighter, cheaper and worked better in damp conditions than Kevlar alone. Kevlar, when used alone, loses about 20% of its effectiveness when wet, so required an expensive waterproofing process. Wool increased friction in a vest with 28–30 layers of fabric, to provide the same level of bullet resistance as 36 layers of Kevlar alone. Events A buyer of Merino wool, Ermenegildo Zegna, has offered awards for Australian wool producers. In 1963, the first Ermenegildo Zegna Perpetual Trophy was presented in Tasmania for growers of "Superfine skirted Merino fleece". In 1980, a national award, the Ermenegildo Zegna Trophy for Extrafine Wool Production, was launched. In 2004, this award became known as the Ermenegildo Zegna Unprotected Wool Trophy. In 1998, an Ermenegildo Zegna Protected Wool Trophy was launched for fleece from sheep coated for around nine months of the year. In 2002, the Ermenegildo Zegna Vellus Aureum Trophy was launched for wool that is 13.9 microns or finer. Wool from Australia, New Zealand, Argentina, and South Africa may enter, and a winner is named from each country. In April 2008, New Zealand won the Ermenegildo Zegna Vellus Aureum Trophy for the first time with a fleece that measured 10.8 microns. This contest awards the winning fleece weight with the same weight in gold as a prize, hence the name. In 2010, an ultrafine, 10-micron fleece, from Windradeen, near Pyramul, New South Wales, won the Ermenegildo Zegna Vellus Aureum International Trophy. Since 2000, Loro Piana has awarded a cup for the world's finest bale of wool that produces just enough fabric for 50 tailor-made suits. The prize is awarded to an Australian or New Zealand wool grower who produces the year's finest bale. The New England Merino Field days which display local studs, wool, and sheep are held during January, in even numbered years around the Walcha, New South Wales district. The Annual Wool Fashion Awards, which showcase the use of Merino wool by fashion designers, are hosted by the city of Armidale, New South Wales, in March each year. This event encourages young and established fashion designers to display their talents. During each May, Armidale hosts the annual New England Wool Expo to display wool fashions, handicrafts, demonstrations, shearing competitions, yard dog trials, and more. In July, the annual Australian Sheep and Wool Show is held in Bendigo, Victoria. This is the largest sheep and wool show in the world, with goats and alpacas, as well as woolcraft competitions and displays, fleece competitions, sheepdog trials, shearing, and wool handling. The largest competition in the world for objectively measured fleeces is the Australian Fleece Competition, which is held annually at Bendigo. In 2008, 475 entries came from all states of Australia, with first and second prizes going to the Northern Tablelands fleeces.	Technology	Fabrics and fibers	null
38436	https://en.wikipedia.org/wiki/Half-reaction	Half-reaction	In chemistry, a half reaction (or half-cell reaction) is either the oxidation or reduction reaction component of a redox reaction. A half reaction is obtained by considering the change in oxidation states of individual substances involved in the redox reaction. Often, the concept of half reactions is used to describe what occurs in an electrochemical cell, such as a Galvanic cell battery. Half reactions can be written to describe both the metal undergoing oxidation (known as the anode) and the metal undergoing reduction (known as the cathode). Half reactions are often used as a method of balancing redox reactions. For oxidation-reduction reactions in acidic conditions, after balancing the atoms and oxidation numbers, one will need to add ions to balance the hydrogen ions in the half reaction. For oxidation-reduction reactions in basic conditions, after balancing the atoms and oxidation numbers, first treat it as an acidic solution and then add ions to balance the ions in the half reactions (which would give ). Example: Zn and Cu Galvanic cell Consider the Galvanic cell shown in the adjacent image: it is constructed with a piece of zinc (Zn) submerged in a solution of zinc sulfate () and a piece of copper (Cu) submerged in a solution of copper(II) sulfate (). The overall reaction is: Zn_{(s)}{} + CuSO4_{(aq)} -> ZnSO4_{(aq)}{} + Cu_{(s)} At the Zn anode, oxidation takes place (the metal loses electrons). This is represented in the following oxidation half reaction (note that the electrons are on the products side): Zn_{(s)} -> Zn^2+ + 2e- At the Cu cathode, reduction takes place (electrons are accepted). This is represented in the following reduction half reaction (note that the electrons are on the reactants side): Cu^2+ + 2e- -> Cu_{(s)} Example: oxidation of magnesium Consider the example burning of magnesium ribbon (Mg). When magnesium burns, it combines with oxygen () from the air to form magnesium oxide (MgO) according to the following equation: 2Mg_{(s)}{} + O2_{(g)} -> 2MgO_{(s)} Magnesium oxide is an ionic compound containing and ions whereas and are elements with no charges. The with zero charge gains a +2 charge going from the reactant side to product side, and the with zero charge gains a –2 charge. This is because when becomes , it loses 2 electrons. Since there are 2 Mg on left side, a total of 4 electrons are lost according to the following oxidation half reaction: 2Mg_{(s)} -> 2Mg^2+ + 4e- On the other hand, was reduced: its oxidation state goes from 0 to -2. Thus, a reduction half reaction can be written for the O2 as it gains 4 electrons: O2_{(g)}{} + 4e- -> 2O^2- The overall reaction is the sum of both half reactions: 2Mg_{(s)}{} + O2_{(g)}{} + 4e- -> 2Mg^2+ + 2O^2- + 4e- When chemical reaction, especially, redox reaction takes place, we do not see the electrons as they appear and disappear during the course of the reaction. What we see is the reactants (starting material) and end products. Due to this, electrons appearing on both sides of the equation are canceled. After canceling, the equation is re-written as 2Mg_{(s)}{} + O2_{(g)} -> 2Mg^2+ + 2O^2- Two ions, positive () and negative () exist on product side and they combine immediately to form a compound magnesium oxide (MgO) due to their opposite charges (electrostatic attraction). In any given oxidation-reduction reaction, there are two half reactions—oxidation half reaction and reduction half reaction. The sum of these two half reactions is the oxidation–reduction reaction. Half-reaction balancing method Consider the reaction below: Cl2 + 2Fe^2+ -> 2Cl- + 2Fe^3+ The two elements involved, iron and chlorine, each change oxidation state; iron from +2 to +3, chlorine from 0 to −1. There are then effectively two half reactions occurring. These changes can be represented in formulas by inserting appropriate electrons into each half reaction: Given two half reactions it is possible, with knowledge of appropriate electrode potentials, to arrive at the complete (original) reaction the same way. The decomposition of a reaction into half reactions is key to understanding a variety of chemical processes. For example, in the above reaction, it can be shown that this is a redox reaction in which Fe is oxidised, and Cl is reduced. Note the transfer of electrons from Fe to Cl. Decomposition is also a way to simplify the balancing of a chemical equation. A chemist can atom balance and charge balance one piece of an equation at a time. For example: becomes is added to and finally becomes It is also possible and sometimes necessary to consider a half reaction in either basic or acidic conditions, as there may be an acidic or basic electrolyte in the redox reaction. Due to this electrolyte it may be more difficult to satisfy the balance of both the atoms and charges. This is done by adding , and/or to either side of the reaction until both atoms and charges are balanced. Consider the half reaction below: PbO2 -> PbO , and can be used to balance the charges and atoms in basic conditions, as long as it is assumed that the reaction is in water. 2e- + H2O + PbO2 -> PbO + 2OH- Again consider the half reaction below: PbO2 -> PbO , and can be used to balance the charges and atoms in acidic conditions, as long as it is assumed that the reaction is in water. 2e- + 2H+ + PbO2 -> PbO + H2O Notice that both sides are both charge balanced and atom balanced. Often there will be both and present in acidic and basic conditions but that the resulting reaction of the two ions will yield water, (shown below): H+ + OH- -> H2O	Physical sciences	Redox reactions	Chemistry
38449	https://en.wikipedia.org/wiki/Affine%20transformation	Affine transformation	In Euclidean geometry, an affine transformation or affinity (from the Latin, affinis, "connected with") is a geometric transformation that preserves lines and parallelism, but not necessarily Euclidean distances and angles. More generally, an affine transformation is an automorphism of an affine space (Euclidean spaces are specific affine spaces), that is, a function which maps an affine space onto itself while preserving both the dimension of any affine subspaces (meaning that it sends points to points, lines to lines, planes to planes, and so on) and the ratios of the lengths of parallel line segments. Consequently, sets of parallel affine subspaces remain parallel after an affine transformation. An affine transformation does not necessarily preserve angles between lines or distances between points, though it does preserve ratios of distances between points lying on a straight line. If is the point set of an affine space, then every affine transformation on can be represented as the composition of a linear transformation on and a translation of . Unlike a purely linear transformation, an affine transformation need not preserve the origin of the affine space. Thus, every linear transformation is affine, but not every affine transformation is linear. Examples of affine transformations include translation, scaling, homothety, similarity, reflection, rotation, hyperbolic rotation, shear mapping, and compositions of them in any combination and sequence. Viewing an affine space as the complement of a hyperplane at infinity of a projective space, the affine transformations are the projective transformations of that projective space that leave the hyperplane at infinity invariant, restricted to the complement of that hyperplane. A generalization of an affine transformation is an affine map (or affine homomorphism or affine mapping) between two (potentially different) affine spaces over the same field . Let and be two affine spaces with and the point sets and and the respective associated vector spaces over the field . A map is an affine map if there exists a linear map such that for all in . Definition Let be an affine space over a field , and be its associated vector space. An affine transformation is a bijection from onto itself that is an affine map; this means that a linear map from to is well defined by the equation here, as usual, the subtraction of two points denotes the free vector from the second point to the first one, and "well-defined" means that implies that If the dimension of is at least two, a semiaffine transformation of is a bijection from onto itself satisfying: For every -dimensional affine subspace of , then is also a -dimensional affine subspace of . If and are parallel affine subspaces of , then and are parallel. These two conditions are satisfied by affine transformations, and express what is precisely meant by the expression that " preserves parallelism". These conditions are not independent as the second follows from the first. Furthermore, if the field has at least three elements, the first condition can be simplified to: is a collineation, that is, it maps lines to lines. Structure By the definition of an affine space, acts on , so that, for every pair in there is associated a point in . We can denote this action by . Here we use the convention that are two interchangeable notations for an element of . By fixing a point in one can define a function by . For any , this function is one-to-one, and so, has an inverse function given by . These functions can be used to turn into a vector space (with respect to the point ) by defining: and This vector space has origin and formally needs to be distinguished from the affine space , but common practice is to denote it by the same symbol and mention that it is a vector space after an origin has been specified. This identification permits points to be viewed as vectors and vice versa. For any linear transformation of , we can define the function by Then is an affine transformation of which leaves the point fixed. It is a linear transformation of , viewed as a vector space with origin . Let be any affine transformation of . Pick a point in and consider the translation of by the vector , denoted by . Translations are affine transformations and the composition of affine transformations is an affine transformation. For this choice of , there exists a unique linear transformation of such that That is, an arbitrary affine transformation of is the composition of a linear transformation of (viewed as a vector space) and a translation of . This representation of affine transformations is often taken as the definition of an affine transformation (with the choice of origin being implicit). Representation As shown above, an affine map is the composition of two functions: a translation and a linear map. Ordinary vector algebra uses matrix multiplication to represent linear maps, and vector addition to represent translations. Formally, in the finite-dimensional case, if the linear map is represented as a multiplication by an invertible matrix and the translation as the addition of a vector , an affine map acting on a vector can be represented as Augmented matrix Using an augmented matrix and an augmented vector, it is possible to represent both the translation and the linear map using a single matrix multiplication. The technique requires that all vectors be augmented with a "1" at the end, and all matrices be augmented with an extra row of zeros at the bottom, an extra column—the translation vector—to the right, and a "1" in the lower right corner. If is a matrix, is equivalent to the following The above-mentioned augmented matrix is called an affine transformation matrix. In the general case, when the last row vector is not restricted to be , the matrix becomes a projective transformation matrix (as it can also be used to perform projective transformations). This representation exhibits the set of all invertible affine transformations as the semidirect product of and . This is a group under the operation of composition of functions, called the affine group. Ordinary matrix-vector multiplication always maps the origin to the origin, and could therefore never represent a translation, in which the origin must necessarily be mapped to some other point. By appending the additional coordinate "1" to every vector, one essentially considers the space to be mapped as a subset of a space with an additional dimension. In that space, the original space occupies the subset in which the additional coordinate is 1. Thus the origin of the original space can be found at . A translation within the original space by means of a linear transformation of the higher-dimensional space is then possible (specifically, a shear transformation). The coordinates in the higher-dimensional space are an example of homogeneous coordinates. If the original space is Euclidean, the higher dimensional space is a real projective space. The advantage of using homogeneous coordinates is that one can combine any number of affine transformations into one by multiplying the respective matrices. This property is used extensively in computer graphics, computer vision and robotics. Example augmented matrix Suppose you have three points that define a non-degenerate triangle in a plane, or four points that define a non-degenerate tetrahedron in 3-dimensional space, or generally points , ..., that define a non-degenerate simplex in -dimensional space. Suppose you have corresponding destination points , ..., , where these new points can lie in a space with any number of dimensions. (Furthermore, the new points need not be distinct from each other and need not form a non-degenerate simplex.) The unique augmented matrix that achieves the affine transformation is Properties Properties preserved An affine transformation preserves: collinearity between points: three or more points which lie on the same line (called collinear points) continue to be collinear after the transformation. parallelism: two or more lines which are parallel, continue to be parallel after the transformation. convexity of sets: a convex set continues to be convex after the transformation. Moreover, the extreme points of the original set are mapped to the extreme points of the transformed set. ratios of lengths of parallel line segments: for distinct parallel segments defined by points and , and , the ratio of and is the same as that of and . barycenters of weighted collections of points. Groups As an affine transformation is invertible, the square matrix appearing in its matrix representation is invertible. The matrix representation of the inverse transformation is thus The invertible affine transformations (of an affine space onto itself) form the affine group, which has the general linear group of degree as subgroup and is itself a subgroup of the general linear group of degree . The similarity transformations form the subgroup where is a scalar times an orthogonal matrix. For example, if the affine transformation acts on the plane and if the determinant of is 1 or −1 then the transformation is an equiareal mapping. Such transformations form a subgroup called the equi-affine group. A transformation that is both equi-affine and a similarity is an isometry of the plane taken with Euclidean distance. Each of these groups has a subgroup of orientation-preserving or positive affine transformations: those where the determinant of is positive. In the last case this is in 3D the group of rigid transformations (proper rotations and pure translations). If there is a fixed point, we can take that as the origin, and the affine transformation reduces to a linear transformation. This may make it easier to classify and understand the transformation. For example, describing a transformation as a rotation by a certain angle with respect to a certain axis may give a clearer idea of the overall behavior of the transformation than describing it as a combination of a translation and a rotation. However, this depends on application and context. Affine maps An affine map between two affine spaces is a map on the points that acts linearly on the vectors (that is, the vectors between points of the space). In symbols, determines a linear transformation such that, for any pair of points : or . We can interpret this definition in a few other ways, as follows. If an origin is chosen, and denotes its image , then this means that for any vector : . If an origin is also chosen, this can be decomposed as an affine transformation that sends , namely , followed by the translation by a vector . The conclusion is that, intuitively, consists of a translation and a linear map. Alternative definition Given two affine spaces and , over the same field, a function is an affine map if and only if for every family of weighted points in such that , we have . In other words, preserves barycenters. History The word "affine" as a mathematical term is defined in connection with tangents to curves in Euler's 1748 Introductio in analysin infinitorum. Felix Klein attributes the term "affine transformation" to Möbius and Gauss. Image transformation In their applications to digital image processing, the affine transformations are analogous to printing on a sheet of rubber and stretching the sheet's edges parallel to the plane. This transform relocates pixels requiring intensity interpolation to approximate the value of moved pixels, bicubic interpolation is the standard for image transformations in image processing applications. Affine transformations scale, rotate, translate, mirror and shear images as shown in the following examples: The affine transforms are applicable to the registration process where two or more images are aligned (registered). An example of image registration is the generation of panoramic images that are the product of multiple images stitched together. Affine warping The affine transform preserves parallel lines. However, the stretching and shearing transformations warp shapes, as the following example shows: This is an example of image warping. However, the affine transformations do not facilitate projection onto a curved surface or radial distortions. In the plane Affine transformations in two real dimensions include: pure translations, scaling in a given direction, with respect to a line in another direction (not necessarily perpendicular), combined with translation that is not purely in the direction of scaling; taking "scaling" in a generalized sense it includes the cases that the scale factor is zero (projection) or negative; the latter includes reflection, and combined with translation it includes glide reflection, rotation combined with a homothety and a translation, shear mapping combined with a homothety and a translation, or squeeze mapping combined with a homothety and a translation. To visualise the general affine transformation of the Euclidean plane, take labelled parallelograms ABCD and A′B′C′D′. Whatever the choices of points, there is an affine transformation T of the plane taking A to A′, and each vertex similarly. Supposing we exclude the degenerate case where ABCD has zero area, there is a unique such affine transformation T. Drawing out a whole grid of parallelograms based on ABCD, the image T(P) of any point P is determined by noting that T(A) = A′, T applied to the line segment AB is A′B′, T applied to the line segment AC is A′C′, and T respects scalar multiples of vectors based at A. [If A, E, F are collinear then the ratio length(AF)/length(AE) is equal to length(A′F′)/length(A′E′).] Geometrically T transforms the grid based on ABCD to that based in A′B′C′D′. Affine transformations do not respect lengths or angles; they multiply area by a constant factor area of A′B′C′D′ / area of ABCD. A given T may either be direct (respect orientation), or indirect (reverse orientation), and this may be determined by its effect on signed areas (as defined, for example, by the cross product of vectors). Examples Over the real numbers The functions with and in and , are precisely the affine transformations of the real line. In plane geometry In , the transformation shown at left is accomplished using the map given by: Transforming the three corner points of the original triangle (in red) gives three new points which form the new triangle (in blue). This transformation skews and translates the original triangle. In fact, all triangles are related to one another by affine transformations. This is also true for all parallelograms, but not for all quadrilaterals.	Mathematics	Non-Euclidean geometry	null
38452	https://en.wikipedia.org/wiki/Oxidation%20state	Oxidation state	In chemistry, the oxidation state, or oxidation number, is the hypothetical charge of an atom if all of its bonds to other atoms were fully ionic. It describes the degree of oxidation (loss of electrons) of an atom in a chemical compound. Conceptually, the oxidation state may be positive, negative or zero. Beside nearly-pure ionic bonding, many covalent bonds exhibit a strong ionicity, making oxidation state a useful predictor of charge. The oxidation state of an atom does not represent the "real" charge on that atom, or any other actual atomic property. This is particularly true of high oxidation states, where the ionization energy required to produce a multiply positive ion is far greater than the energies available in chemical reactions. Additionally, the oxidation states of atoms in a given compound may vary depending on the choice of electronegativity scale used in their calculation. Thus, the oxidation state of an atom in a compound is purely a formalism. It is nevertheless important in understanding the nomenclature conventions of inorganic compounds. Also, several observations regarding chemical reactions may be explained at a basic level in terms of oxidation states. Oxidation states are typically represented by integers which may be positive, zero, or negative. In some cases, the average oxidation state of an element is a fraction, such as for iron in magnetite (see below). The highest known oxidation state is reported to be +9, displayed by iridium in the tetroxoiridium(IX) cation (). It is predicted that even a +10 oxidation state may be achieved by platinum in tetroxoplatinum(X), . The lowest oxidation state is −5, as for boron in and gallium in pentamagnesium digallide (). In Stock nomenclature, which is commonly used for inorganic compounds, the oxidation state is represented by a Roman numeral placed after the element name inside parentheses or as a superscript after the element symbol, e.g. Iron(III) oxide. The term oxidation was first used by Antoine Lavoisier to signify the reaction of a substance with oxygen. Much later, it was realized that the substance, upon being oxidized, loses electrons, and the meaning was extended to include other reactions in which electrons are lost, regardless of whether oxygen was involved. The increase in the oxidation state of an atom, through a chemical reaction, is known as oxidation; a decrease in oxidation state is known as a reduction. Such reactions involve the formal transfer of electrons: a net gain in electrons being a reduction, and a net loss of electrons being oxidation. For pure elements, the oxidation state is zero. Overview Oxidation numbers are assigned to elements in a molecule such that the overall sum is zero in a neutral molecule. The number indicates the degree of oxidation of each element caused by molecular bonding. In ionic compounds, the oxidation numbers are the same as the element's ionic charge. Thus for KCl, potassium is assigned +1 and chlorine is assigned -1. The complete set of rules for assigning oxidation numbers are discussed in the following sections. Oxidation numbers are fundamental to the chemical nomenclature of ionic compounds. For example, Cu compounds with Cu oxidation state +2 are call cupric and those with state +1 are cuprous. The oxidation numbers of elements allow predictions of chemical formula and reactions, especially oxidation-reduction reactions. The oxidation numbers of the most stable chemical compounds follow trends in the periodic table. IUPAC definition International Union of Pure and Applied Chemistry (IUPAC) has published a "Comprehensive definition of oxidation state (IUPAC Recommendations 2016)". It is a distillation of an IUPAC technical report "Toward a comprehensive definition of oxidation state" from 2014. The current IUPAC Gold Book definition of oxidation state is: and the term oxidation number is nearly synonymous. The ionic approximation means extrapolating bonds to ionic. Several criteria were considered for the ionic approximation: Extrapolation of the bond's polarity; Assignment of electrons according to the atom's contribution to the bonding Molecular orbital (MO) or the electron's allegiance in a LCAO–MO model. In a bond between two different elements, the bond's electrons are assigned to its main atomic contributor typically of higher electronegativity; in a bond between two atoms of the same element, the electrons are divided equally. Most electronegativity scales depend on the atom's bonding state, which makes the assignment of the oxidation state a somewhat circular argument. For example, some scales may turn out unusual oxidation states, such as −6 for platinum in , for Pauling and Mulliken scales. The dipole moments would, sometimes, also turn out abnormal oxidation numbers, such as in CO and NO, which are oriented with their positive end towards oxygen. Therefore, this leaves the atom's contribution to the bonding MO, the atomic-orbital energy, and from quantum-chemical calculations of charges, as the only viable criteria with cogent values for ionic approximation. However, for a simple estimate for the ionic approximation, we can use Allen electronegativities, as only that electronegativity scale is truly independent of the oxidation state, as it relates to the average valence‐electron energy of the free atom: Determination While introductory levels of chemistry teaching use postulated oxidation states, the IUPAC recommendation and the Gold Book entry list two entirely general algorithms for the calculation of the oxidation states of elements in chemical compounds. Simple approach without bonding considerations Introductory chemistry uses postulates: the oxidation state for an element in a chemical formula is calculated from the overall charge and postulated oxidation states for all the other atoms. A simple example is based on two postulates, OS = +1 for hydrogen OS = −2 for oxygen where OS stands for oxidation state. This approach yields correct oxidation states in oxides and hydroxides of any single element, and in acids such as sulfuric acid () or dichromic acid (). Its coverage can be extended either by a list of exceptions or by assigning priority to the postulates. The latter works for hydrogen peroxide () where the priority of rule 1 leaves both oxygens with oxidation state −1. Additional postulates and their ranking may expand the range of compounds to fit a textbook's scope. As an example, one postulatory algorithm from many possible; in a sequence of decreasing priority: An element in a free form has OS = 0. In a compound or ion, the sum of the oxidation states equals the total charge of the compound or ion. Fluorine in compounds has OS = −1; this extends to chlorine and bromine only when not bonded to a lighter halogen, oxygen or nitrogen. Group 1 and group 2 metals in compounds have OS = +1 and +2, respectively. Hydrogen has OS = +1 but adopts −1 when bonded as a hydride to metals or metalloids. Oxygen in compounds has OS = −2 but only when not bonded to oxygen (e.g. in peroxides) or fluorine. This set of postulates covers oxidation states of fluorides, chlorides, bromides, oxides, hydroxides, and hydrides of any single element. It covers all oxoacids of any central atom (and all their fluoro-, chloro-, and bromo-relatives), as well as salts of such acids with group 1 and 2 metals. It also covers iodides, sulfides, and similar simple salts of these metals. Algorithm of assigning bonds This algorithm is performed on a Lewis structure (a diagram that shows all valence electrons). Oxidation state equals the charge of an atom after each of its heteronuclear bonds has been assigned to the more electronegative partner of the bond (except when that partner is a reversibly bonded Lewis-acid ligand) and homonuclear bonds have been divided equally: where each "—" represents an electron pair (either shared between two atoms or solely on one atom), and "OS" is the oxidation state as a numerical variable. After the electrons have been assigned according to the vertical red lines on the formula, the total number of valence electrons that now "belong" to each atom is subtracted from the number of valence electrons of the neutral atom (such as 5 for nitrogen in group 15) to yield that atom's oxidation state. This example shows the importance of describing the bonding. Its summary formula, , corresponds to two structural isomers; the peroxynitrous acid in the above figure and the more stable nitric acid. With the formula , the simple approach without bonding considerations yields −2 for all three oxygens and +5 for nitrogen, which is correct for nitric acid. For the peroxynitrous acid, however, both oxygens in the O–O bond have OS = −1, and the nitrogen has OS = +3, which requires a structure to understand. Organic compounds are treated in a similar manner; exemplified here on functional groups occurring in between methane () and carbon dioxide (): Analogously for transition-metal compounds; on the left has a total of 36 valence electrons (18 pairs to be distributed), and hexacarbonylchromium () on the right has 66 valence electrons (33 pairs): A key step is drawing the Lewis structure of the molecule (neutral, cationic, anionic): Atom symbols are arranged so that pairs of atoms can be joined by single two-electron bonds as in the molecule (a sort of "skeletal" structure), and the remaining valence electrons are distributed such that sp atoms obtain an octet (duet for hydrogen) with a priority that increases in proportion with electronegativity. In some cases, this leads to alternative formulae that differ in bond orders (the full set of which is called the resonance formulas). Consider the sulfate anion () with 32 valence electrons; 24 from oxygens, 6 from sulfur, 2 of the anion charge obtained from the implied cation. The bond orders to the terminal oxygens do not affect the oxidation state so long as the oxygens have octets. Already the skeletal structure, top left, yields the correct oxidation states, as does the Lewis structure, top right (one of the resonance formulas): The bond-order formula at the bottom is closest to the reality of four equivalent oxygens each having a total bond order of 2. That total includes the bond of order to the implied cation and follows the 8 − N rule requiring that the main-group atom's bond-order total equals 8 − N valence electrons of the neutral atom, enforced with a priority that proportionately increases with electronegativity. This algorithm works equally for molecular cations composed of several atoms. An example is the ammonium cation of 8 valence electrons (5 from nitrogen, 4 from hydrogens, minus 1 electron for the cation's positive charge): Drawing Lewis structures with electron pairs as dashes emphasizes the essential equivalence of bond pairs and lone pairs when counting electrons and moving bonds onto atoms. Structures drawn with electron dot pairs are of course identical in every way: The algorithm's caveat The algorithm contains a caveat, which concerns rare cases of transition-metal complexes with a type of ligand that is reversibly bonded as a Lewis acid (as an acceptor of the electron pair from the transition metal); termed a "Z-type" ligand in Green's covalent bond classification method. The caveat originates from the simplifying use of electronegativity instead of the MO-based electron allegiance to decide the ionic sign. One early example is the complex with sulfur dioxide () as the reversibly-bonded acceptor ligand (released upon heating). The Rh−S bond is therefore extrapolated ionic against Allen electronegativities of rhodium and sulfur, yielding oxidation state +1 for rhodium: Algorithm of summing bond orders This algorithm works on Lewis structures and bond graphs of extended (non-molecular) solids: Applied to a Lewis structure An example of a Lewis structure with no formal charge, illustrates that, in this algorithm, homonuclear bonds are simply ignored (the bond orders are in blue). Carbon monoxide exemplifies a Lewis structure with formal charges: To obtain the oxidation states, the formal charges are summed with the bond-order value taken positively at the carbon and negatively at the oxygen. Applied to molecular ions, this algorithm considers the actual location of the formal (ionic) charge, as drawn in the Lewis structure. As an example, summing bond orders in the ammonium cation yields −4 at the nitrogen of formal charge +1, with the two numbers adding to the oxidation state of −3: The sum of oxidation states in the ion equals its charge (as it equals zero for a neutral molecule). Also in anions, the formal (ionic) charges have to be considered when nonzero. For sulfate this is exemplified with the skeletal or Lewis structures (top), compared with the bond-order formula of all oxygens equivalent and fulfilling the octet and 8 − N rules (bottom): Applied to bond graph A bond graph in solid-state chemistry is a chemical formula of an extended structure, in which direct bonding connectivities are shown. An example is the perovskite, the unit cell of which is drawn on the left and the bond graph (with added numerical values) on the right: We see that the oxygen atom bonds to the six nearest rubidium cations, each of which has 4 bonds to the auride anion. The bond graph summarizes these connectivities. The bond orders (also called bond valences) sum up to oxidation states according to the attached sign of the bond's ionic approximation (there are no formal charges in bond graphs). Determination of oxidation states from a bond graph can be illustrated on ilmenite, . We may ask whether the mineral contains and , or and . Its crystal structure has each metal atom bonded to six oxygens and each of the equivalent oxygens to two irons and two titaniums, as in the bond graph below. Experimental data show that three metal-oxygen bonds in the octahedron are short and three are long (the metals are off-center). The bond orders (valences), obtained from the bond lengths by the bond valence method, sum up to 2.01 at Fe and 3.99 at Ti; which can be rounded off to oxidation states +2 and +4, respectively: Balancing redox Oxidation states can be useful for balancing chemical equations for oxidation-reduction (or redox) reactions, because the changes in the oxidized atoms have to be balanced by the changes in the reduced atoms. For example, in the reaction of acetaldehyde with Tollens' reagent to form acetic acid (shown below), the carbonyl carbon atom changes its oxidation state from +1 to +3 (loses two electrons). This oxidation is balanced by reducing two cations to (gaining two electrons in total). An inorganic example is the Bettendorf reaction using tin dichloride () to prove the presence of arsenite ions in a concentrated HCl extract. When arsenic(III) is present, a brown coloration appears forming a dark precipitate of arsenic, according to the following simplified reaction: Here three tin atoms are oxidized from oxidation state +2 to +4, yielding six electrons that reduce two arsenic atoms from oxidation state +3 to 0. The simple one-line balancing goes as follows: the two redox couples are written down as they react; One tin is oxidized from oxidation state +2 to +4, a two-electron step, hence 2 is written in front of the two arsenic partners. One arsenic is reduced from +3 to 0, a three-electron step, hence 3 goes in front of the two tin partners. An alternative three-line procedure is to write separately the half-reactions for oxidation and reduction, each balanced with electrons, and then to sum them up such that the electrons cross out. In general, these redox balances (the one-line balance or each half-reaction) need to be checked for the ionic and electron charge sums on both sides of the equation being indeed equal. If they are not equal, suitable ions are added to balance the charges and the non-redox elemental balance. Appearances Nominal oxidation states A nominal oxidation state is a general term with two different definitions: Electrochemical oxidation state represents a molecule or ion in the Latimer diagram or Frost diagram for its redox-active element. An example is the Latimer diagram for sulfur at pH 0 where the electrochemical oxidation state +2 for sulfur puts between S and H2SO3: Systematic oxidation state is chosen from close alternatives as a pedagogical description. An example is the oxidation state of phosphorus in H3PO3 (structurally diprotic HPO(OH)2) taken nominally as +3, while Allen electronegativities of phosphorus and hydrogen suggest +5 by a narrow margin that makes the two alternatives almost equivalent: Both alternative oxidation numbers for phosphorus make chemical sense, depending on which chemical property or reaction is emphasized. By contrast, a calculated alternative, such as the average (+4) does not. Ambiguous oxidation states Lewis formulae are rule-based approximations of chemical reality, as are Allen electronegativities. Still, oxidation states may seem ambiguous when their determination is not straightforward. If only an experiment can determine the oxidation state, the rule-based determination is ambiguous (insufficient). There are also truly dichotomous values that are decided arbitrarily. Oxidation-state determination from resonance formulas Seemingly ambiguous oxidation states are derived from a set of resonance formulas of equal weights for a molecule having heteronuclear bonds where the atom connectivity does not correspond to the number of two-electron bonds dictated by the 8 − N rule. An example is S2N2 where four resonance formulas featuring one S=N double bond have oxidation states +2 and +4 for the two sulfur atoms, which average to +3 because the two sulfur atoms are equivalent in this square-shaped molecule. A physical measurement is needed to determine oxidation state when a non-innocent ligand is present, of hidden or unexpected redox properties that could otherwise be assigned to the central atom. An example is the nickel dithiolate complex, . when the redox ambiguity of a central atom and ligand yields dichotomous oxidation states of close stability, thermally induced tautomerism may result, as exemplified by manganese catecholate, . Assignment of such oxidation states requires spectroscopic, magnetic or structural data. when the bond order has to be ascertained along with an isolated tandem of a heteronuclear and a homonuclear bond. An example is thiosulfate having two possible oxidation states (bond orders are in blue and formal charges in green): The S–S distance measurement in thiosulfate is needed to reveal that this bond order is very close to 1, as in the formula on the left. Ambiguous/arbitrary oxidation states when the electronegativity difference between two bonded atoms is very small (as in H3PO3). Two almost equivalent pairs of oxidation states, arbitrarily chosen, are obtained for these atoms. when an electronegative p-block atom forms solely homonuclear bonds, the number of which differs from the number of two-electron bonds suggested by rules. Examples are homonuclear finite chains like (the central nitrogen connects two atoms with four two-electron bonds while only three two-electron bonds are required by the 8 − N rule) or (the central iodine connects two atoms with two two-electron bonds while only one two-electron bond fulfills the 8 − N rule). A sensible approach is to distribute the ionic charge over the two outer atoms. Such a placement of charges in a polysulfide (where all inner sulfurs form two bonds, fulfilling the 8 − N rule) follows already from its Lewis structure. when the isolated tandem of a heteronuclear and a homonuclear bond leads to a bonding compromise in between two Lewis structures of limiting bond orders. An example is N2O: The typical oxidation state of nitrogen in N2O is +1, which also obtains for both nitrogens by a molecular orbital approach. The formal charges on the right comply with electronegativities, which implies an added ionic bonding contribution. Indeed, the estimated N−N and N−O bond orders are 2.76 and 1.9, respectively, approaching the formula of integer bond orders that would include the ionic contribution explicitly as a bond (in green): Conversely, formal charges against electronegativities in a Lewis structure decrease the bond order of the corresponding bond. An example is carbon monoxide with a bond-order estimate of 2.6. Fractional oxidation states Fractional oxidation states are often used to represent the average oxidation state of several atoms of the same element in a structure. For example, the formula of magnetite is , implying an average oxidation state for iron of +. However, this average value may not be representative if the atoms are not equivalent. In a crystal below , two-thirds of the cations are and one-third are , and the formula may be more clearly represented as FeO·. Likewise, propane, , has been described as having a carbon oxidation state of −. Again, this is an average value since the structure of the molecule is , with the first and third carbon atoms each having an oxidation state of −3 and the central one −2. An example with true fractional oxidation states for equivalent atoms is potassium superoxide, . The diatomic superoxide ion has an overall charge of −1, so each of its two equivalent oxygen atoms is assigned an oxidation state of −. This ion can be described as a resonance hybrid of two Lewis structures, where each oxygen has an oxidation state of 0 in one structure and −1 in the other. For the cyclopentadienyl anion , the oxidation state of C is −1 + − = −. The −1 occurs because each carbon is bonded to one hydrogen atom (a less electronegative element), and the − because the total ionic charge of −1 is divided among five equivalent carbons. Again this can be described as a resonance hybrid of five equivalent structures, each having four carbons with oxidation state −1 and one with −2. {\| class="wikitable" \|+ Examples of fractional oxidation states for carbon \|- ! Oxidation state !! Example species \|- \| − \|\| \|- \| − \|\| \|- \| + \|\| \|} Finally, fractional oxidation numbers are not used in the chemical nomenclature. For example the red lead is represented as lead(II,IV) oxide, showing the oxidation states of the two nonequivalent lead atoms. Elements with multiple oxidation states Most elements have more than one possible oxidation state. For example, carbon has nine possible integer oxidation states from −4 to +4: {\| class="wikitable" \|+ Integer oxidation states of carbon \|- ! Oxidation state !! Example compound \|- \| −4 \|\| \|- \| −3 \|\| \|- \| −2 \|\| , \|- \| −1 \|\| , , \|- \| 0 \|\| , \|- \| +1 \|\| , \|- \| +2 \|\| , \|- \| +3 \|\| , \|- \| +4 \|\| , \|} Oxidation state in metals Many compounds with luster and electrical conductivity maintain a simple stoichiometric formula, such as the golden TiO, blue-black RuO2 or coppery ReO3, all of obvious oxidation state. Ultimately, assigning the free metallic electrons to one of the bonded atoms is not comprehensive and can yield unusual oxidation states. Examples are the LiPb and ordered alloys, the composition and structure of which are largely determined by atomic size and packing factors. Should oxidation state be needed for redox balancing, it is best set to 0 for all atoms of such an alloy. List of oxidation states of the elements This is a list of known oxidation states of the chemical elements, excluding nonintegral values. The most common states appear in bold. The table is based on that of Greenwood and Earnshaw, with additions noted. Every element exists in oxidation state 0 when it is the pure non-ionized element in any phase, whether monatomic or polyatomic allotrope. The column for oxidation state 0 only shows elements known to exist in oxidation state 0 in compounds. Early forms (octet rule) A figure with a similar format was used by Irving Langmuir in 1919 in one of the early papers about the octet rule. The periodicity of the oxidation states was one of the pieces of evidence that led Langmuir to adopt the rule. Use in nomenclature The oxidation state in compound naming for transition metals and lanthanides and actinides is placed either as a right superscript to the element symbol in a chemical formula, such as FeIII or in parentheses after the name of the element in chemical names, such as iron(III). For example, is named iron(III) sulfate and its formula can be shown as Fe. This is because a sulfate ion has a charge of −2, so each iron atom takes a charge of +3. History of the oxidation state concept Early days Oxidation itself was first studied by Antoine Lavoisier, who defined it as the result of reactions with oxygen (hence the name). The term has since been generalized to imply a formal loss of electrons. Oxidation states, called oxidation grades by Friedrich Wöhler in 1835, were one of the intellectual stepping stones that Dmitri Mendeleev used to derive the periodic table. William B. Jensen gives an overview of the history up to 1938. Use in nomenclature When it was realized that some metals form two different binary compounds with the same nonmetal, the two compounds were often distinguished by using the ending -ic for the higher metal oxidation state and the ending -ous for the lower. For example, FeCl3 is ferric chloride and FeCl2 is ferrous chloride. This system is not very satisfactory (although sometimes still used) because different metals have different oxidation states which have to be learned: ferric and ferrous are +3 and +2 respectively, but cupric and cuprous are +2 and +1, and stannic and stannous are +4 and +2. Also, there was no allowance for metals with more than two oxidation states, such as vanadium with oxidation states +2, +3, +4, and +5. This system has been largely replaced by one suggested by Alfred Stock in 1919 and adopted by IUPAC in 1940. Thus, FeCl2 was written as iron(II) chloride rather than ferrous chloride. The Roman numeral II at the central atom came to be called the "Stock number" (now an obsolete term), and its value was obtained as a charge at the central atom after removing its ligands along with the electron pairs they shared with it. Development towards the current concept The term "oxidation state" in English chemical literature was popularized by Wendell Mitchell Latimer in his 1938 book about electrochemical potentials. He used it for the value (synonymous with the German term Wertigkeit) previously termed "valence", "polar valence" or "polar number" in English, or "oxidation stage" or indeed the "state of oxidation". Since 1938, the term "oxidation state" has been connected with electrochemical potentials and electrons exchanged in redox couples participating in redox reactions. By 1948, IUPAC used the 1940 nomenclature rules with the term "oxidation state", instead of the original valency. In 1948 Linus Pauling proposed that oxidation number could be determined by extrapolating bonds to being completely ionic in the direction of electronegativity. A full acceptance of this suggestion was complicated by the fact that the Pauling electronegativities as such depend on the oxidation state and that they may lead to unusual values of oxidation states for some transition metals. In 1990 IUPAC resorted to a postulatory (rule-based) method to determine the oxidation state. This was complemented by the synonymous term oxidation number as a descendant of the Stock number introduced in 1940 into the nomenclature. However, the terminology using "ligands" gave the impression that oxidation number might be something specific to coordination complexes. This situation and the lack of a real single definition generated numerous debates about the meaning of oxidation state, suggestions about methods to obtain it and definitions of it. To resolve the issue, an IUPAC project (2008-040-1-200) was started in 2008 on the "Comprehensive Definition of Oxidation State", and was concluded by two reports and by the revised entries "Oxidation State" and "Oxidation Number" in the IUPAC Gold Book. The outcomes were a single definition of oxidation state and two algorithms to calculate it in molecular and extended-solid compounds, guided by Allen electronegativities that are independent of oxidation state.	Physical sciences	Redox reactions	Chemistry
38454	https://en.wikipedia.org/wiki/Gravitational%20constant	Gravitational constant	The gravitational constant is an empirical physical constant involved in the calculation of gravitational effects in Sir Isaac Newton's law of universal gravitation and in Albert Einstein's theory of general relativity. It is also known as the universal gravitational constant, the Newtonian constant of gravitation, or the Cavendish gravitational constant, denoted by the capital letter . In Newton's law, it is the proportionality constant connecting the gravitational force between two bodies with the product of their masses and the inverse square of their distance. In the Einstein field equations, it quantifies the relation between the geometry of spacetime and the energy–momentum tensor (also referred to as the stress–energy tensor). The measured value of the constant is known with some certainty to four significant digits. In SI units, its value is approximately The modern notation of Newton's law involving was introduced in the 1890s by C. V. Boys. The first implicit measurement with an accuracy within about 1% is attributed to Henry Cavendish in a 1798 experiment. Definition According to Newton's law of universal gravitation, the magnitude of the attractive force () between two bodies each with a spherically symmetric density distribution is directly proportional to the product of their masses, and , and inversely proportional to the square of the distance, , directed along the line connecting their centres of mass: The constant of proportionality, , in this non-relativistic formulation is the gravitational constant. Colloquially, the gravitational constant is also called "Big G", distinct from "small g" (), which is the local gravitational field of Earth (also referred to as free-fall acceleration). Where is the mass of the Earth and is the radius of the Earth, the two quantities are related by: The gravitational constant appears in the Einstein field equations of general relativity, where is the Einstein tensor (not the gravitational constant despite the use of ), is the cosmological constant, is the metric tensor, is the stress–energy tensor, and is the Einstein gravitational constant, a constant originally introduced by Einstein that is directly related to the Newtonian constant of gravitation: Value and uncertainty The gravitational constant is a physical constant that is difficult to measure with high accuracy. This is because the gravitational force is an extremely weak force as compared to other fundamental forces at the laboratory scale. In SI units, the CODATA-recommended value of the gravitational constant is: = The relative standard uncertainty is . Natural units Due to its use as a defining constant in some systems of natural units, particularly geometrized unit systems such as Planck units and Stoney units, the value of the gravitational constant will generally have a numeric value of 1 or a value close to it when expressed in terms of those units. Due to the significant uncertainty in the measured value of G in terms of other known fundamental constants, a similar level of uncertainty will show up in the value of many quantities when expressed in such a unit system. Orbital mechanics In astrophysics, it is convenient to measure distances in parsecs (pc), velocities in kilometres per second (km/s) and masses in solar units . In these units, the gravitational constant is: For situations where tides are important, the relevant length scales are solar radii rather than parsecs. In these units, the gravitational constant is: In orbital mechanics, the period of an object in circular orbit around a spherical object obeys where is the volume inside the radius of the orbit, and is the total mass of the two objects. It follows that This way of expressing shows the relationship between the average density of a planet and the period of a satellite orbiting just above its surface. For elliptical orbits, applying Kepler's 3rd law, expressed in units characteristic of Earth's orbit: where distance is measured in terms of the semi-major axis of Earth's orbit (the astronomical unit, AU), time in years, and mass in the total mass of the orbiting system (). The above equation is exact only within the approximation of the Earth's orbit around the Sun as a two-body problem in Newtonian mechanics, the measured quantities contain corrections from the perturbations from other bodies in the solar system and from general relativity. From 1964 until 2012, however, it was used as the definition of the astronomical unit and thus held by definition: Since 2012, the AU is defined as exactly, and the equation can no longer be taken as holding precisely. The quantity —the product of the gravitational constant and the mass of a given astronomical body such as the Sun or Earth—is known as the standard gravitational parameter (also denoted ). The standard gravitational parameter appears as above in Newton's law of universal gravitation, as well as in formulas for the deflection of light caused by gravitational lensing, in Kepler's laws of planetary motion, and in the formula for escape velocity. This quantity gives a convenient simplification of various gravity-related formulas. The product is known much more accurately than either factor is. Calculations in celestial mechanics can also be carried out using the units of solar masses, mean solar days and astronomical units rather than standard SI units. For this purpose, the Gaussian gravitational constant was historically in widespread use, , expressing the mean angular velocity of the Sun–Earth system. The use of this constant, and the implied definition of the astronomical unit discussed above, has been deprecated by the IAU since 2012. History of measurement Early history The existence of the constant is implied in Newton's law of universal gravitation as published in the 1680s (although its notation as dates to the 1890s), but is not calculated in his Philosophiæ Naturalis Principia Mathematica where it postulates the inverse-square law of gravitation. In the Principia, Newton considered the possibility of measuring gravity's strength by measuring the deflection of a pendulum in the vicinity of a large hill, but thought that the effect would be too small to be measurable. Nevertheless, he had the opportunity to estimate the order of magnitude of the constant when he surmised that "the mean density of the earth might be five or six times as great as the density of water", which is equivalent to a gravitational constant of the order: ≈ A measurement was attempted in 1738 by Pierre Bouguer and Charles Marie de La Condamine in their "Peruvian expedition". Bouguer downplayed the significance of their results in 1740, suggesting that the experiment had at least proved that the Earth could not be a hollow shell, as some thinkers of the day, including Edmond Halley, had suggested. The Schiehallion experiment, proposed in 1772 and completed in 1776, was the first successful measurement of the mean density of the Earth, and thus indirectly of the gravitational constant. The result reported by Charles Hutton (1778) suggested a density of ( times the density of water), about 20% below the modern value. This immediately led to estimates on the densities and masses of the Sun, Moon and planets, sent by Hutton to Jérôme Lalande for inclusion in his planetary tables. As discussed above, establishing the average density of Earth is equivalent to measuring the gravitational constant, given Earth's mean radius and the mean gravitational acceleration at Earth's surface, by setting Based on this, Hutton's 1778 result is equivalent to . The first direct measurement of gravitational attraction between two bodies in the laboratory was performed in 1798, seventy-one years after Newton's death, by Henry Cavendish. He determined a value for implicitly, using a torsion balance invented by the geologist Rev. John Michell (1753). He used a horizontal torsion beam with lead balls whose inertia (in relation to the torsion constant) he could tell by timing the beam's oscillation. Their faint attraction to other balls placed alongside the beam was detectable by the deflection it caused. In spite of the experimental design being due to Michell, the experiment is now known as the Cavendish experiment for its first successful execution by Cavendish. Cavendish's stated aim was the "weighing of Earth", that is, determining the average density of Earth and the Earth's mass. His result, , corresponds to value of . It is surprisingly accurate, about 1% above the modern value (comparable to the claimed relative standard uncertainty of 0.6%). 19th century The accuracy of the measured value of has increased only modestly since the original Cavendish experiment. is quite difficult to measure because gravity is much weaker than other fundamental forces, and an experimental apparatus cannot be separated from the gravitational influence of other bodies. Measurements with pendulums were made by Francesco Carlini (1821, ), Edward Sabine (1827, ), Carlo Ignazio Giulio (1841, ) and George Biddell Airy (1854, ). Cavendish's experiment was first repeated by Ferdinand Reich (1838, 1842, 1853), who found a value of , which is actually worse than Cavendish's result, differing from the modern value by 1.5%. Cornu and Baille (1873), found . Cavendish's experiment proved to result in more reliable measurements than pendulum experiments of the "Schiehallion" (deflection) type or "Peruvian" (period as a function of altitude) type. Pendulum experiments still continued to be performed, by Robert von Sterneck (1883, results between 5.0 and ) and Thomas Corwin Mendenhall (1880, ). Cavendish's result was first improved upon by John Henry Poynting (1891), who published a value of , differing from the modern value by 0.2%, but compatible with the modern value within the cited relative standard uncertainty of 0.55%. In addition to Poynting, measurements were made by C. V. Boys (1895) and Carl Braun (1897), with compatible results suggesting = . The modern notation involving the constant was introduced by Boys in 1894 and becomes standard by the end of the 1890s, with values usually cited in the cgs system. Richarz and Krigar-Menzel (1898) attempted a repetition of the Cavendish experiment using 100,000 kg of lead for the attracting mass. The precision of their result of was, however, of the same order of magnitude as the other results at the time. Arthur Stanley Mackenzie in The Laws of Gravitation (1899) reviews the work done in the 19th century. Poynting is the author of the article "Gravitation" in the Encyclopædia Britannica Eleventh Edition (1911). Here, he cites a value of = with a relative uncertainty of 0.2%. Modern value Paul R. Heyl (1930) published the value of (relative uncertainty 0.1%), improved to (relative uncertainty 0.045% = 450 ppm) in 1942. However, Heyl used the statistical spread as his standard deviation, and he admitted himself that measurements using the same material yielded very similar results while measurements using different materials yielded vastly different results. He spent the next 12 years after his 1930 paper to do more precise measurements, hoping that the composition-dependent effect would go away, but it did not, as he noted in his final paper from the year 1942. Published values of derived from high-precision measurements since the 1950s have remained compatible with Heyl (1930), but within the relative uncertainty of about 0.1% (or 1000 ppm) have varied rather broadly, and it is not entirely clear if the uncertainty has been reduced at all since the 1942 measurement. Some measurements published in the 1980s to 2000s were, in fact, mutually exclusive. Establishing a standard value for with a relative standard uncertainty better than 0.1% has therefore remained rather speculative. By 1969, the value recommended by the National Institute of Standards and Technology (NIST) was cited with a relative standard uncertainty of 0.046% (460 ppm), lowered to 0.012% (120 ppm) by 1986. But the continued publication of conflicting measurements led NIST to considerably increase the standard uncertainty in the 1998 recommended value, by a factor of 12, to a standard uncertainty of 0.15%, larger than the one given by Heyl (1930). The uncertainty was again lowered in 2002 and 2006, but once again raised, by a more conservative 20%, in 2010, matching the relative standard uncertainty of 120 ppm published in 1986. For the 2014 update, CODATA reduced the uncertainty to 46 ppm, less than half the 2010 value, and one order of magnitude below the 1969 recommendation. The following table shows the NIST recommended values published since 1969: In the January 2007 issue of Science, Fixler et al. described a measurement of the gravitational constant by a new technique, atom interferometry, reporting a value of , 0.28% (2800 ppm) higher than the 2006 CODATA value. An improved cold atom measurement by Rosi et al. was published in 2014 of . Although much closer to the accepted value (suggesting that the Fixler et al. measurement was erroneous), this result was 325 ppm below the recommended 2014 CODATA value, with non-overlapping standard uncertainty intervals. As of 2018, efforts to re-evaluate the conflicting results of measurements are underway, coordinated by NIST, notably a repetition of the experiments reported by Quinn et al. (2013). In August 2018, a Chinese research group announced new measurements based on torsion balances, and based on two different methods. These are claimed as the most accurate measurements ever made, with standard uncertainties cited as low as 12 ppm. The difference of 2.7σ between the two results suggests there could be sources of error unaccounted for. Constancy Analysis of observations of 580 type Ia supernovae shows that the gravitational constant has varied by less than one part in ten billion per year over the last nine billion years.	Physical sciences	Physical constants	Physics
38468	https://en.wikipedia.org/wiki/Psilocybin	Psilocybin	Psilocybin, also known as 4-phosphoryloxy-N,N-dimethyltryptamine (4-PO-DMT), and formerly sold under the brand name Indocybin, is a naturally occurring psychedelic prodrug compound produced by more than 200 species of fungi. Psilocybin is itself biologically inactive but is quickly converted by the body to psilocin, which has mind-altering effects similar, in some aspects, to those of other classical psychedelics. Effects include euphoria, hallucinations, changes in perception, a distorted sense of time, and perceived spiritual experiences. It can also cause adverse reactions such as nausea and panic attacks. Imagery in cave paintings and rock art of modern-day Algeria and Spain suggests that human use of psilocybin mushrooms predates recorded history. In Mesoamerica, the mushrooms had long been consumed in spiritual and divinatory ceremonies before Spanish chroniclers first documented their use in the 16th century. In 1958, the Swiss chemist Albert Hofmann isolated psilocybin and psilocin from the mushroom Psilocybe mexicana. His employer, Sandoz, marketed and sold pure psilocybin to physicians and clinicians worldwide for use in psychedelic therapy. Increasingly restrictive drug laws of the 1960s and the 1970s curbed scientific research into the effects of psilocybin and other hallucinogens, but its popularity as an entheogen (spirituality-enhancing agent) grew in the next decade, owing largely to the increased availability of information on how to cultivate psilocybin mushrooms. The intensity and duration of psilocybin's effects vary, depending on species or cultivar of mushrooms, dosage, individual physiology, and set and setting, as shown in experiments led by Timothy Leary at Harvard University in the early 1960s. Once ingested, psilocybin is rapidly metabolized to psilocin, which then acts on serotonin receptors in the brain. Psilocybin's mind-altering effects typically last two to six hours, although to people under the influence of psilocybin, they may seem to last much longer, since the drug can distort the perception of time. Possession of psilocybin-containing mushrooms has been outlawed in most countries, and psilocybin has been classified as a Schedule I controlled substance under the 1971 United Nations Convention on Psychotropic Substances. History Early There is evidence to suggest that psychoactive mushrooms have been used by humans in religious ceremonies for thousands of years. The Tassili Mushroom Figure was discovered in Tassili, Algeria, and is believed to depict psychedelic mushrooms and the transformation of the user under their influence. The paintings are said to date back to 9000-7000 BC. 6,000-year-old pictographs discovered near the Spanish town of Villar del Humo illustrate several mushrooms that have been tentatively identified as Psilocybe hispanica, a hallucinogenic species native to the area. Some scholars have also interpreted archaeological artifacts from Mexico and the so-called Mayan "mushroom stones" of Guatemala as evidence of ritual and ceremonial use of psychoactive mushrooms in the Mayan and Aztec cultures of Mesoamerica. In Nahuatl, the language of the Aztecs, the mushrooms were called teonanácatl—literally "divine mushroom": the agglutinative form of teō(tl) ("god", "sacred") and nanācatl ("mushroom") in Nahuatl. After Spanish explorers of the New World arrived in the 16th century, chroniclers reported the use of mushrooms by the natives for ceremonial and religious purposes. According to the Dominican friar Diego Durán in The History of the Indies of New Spain (published c. 1581), mushrooms were eaten in festivities conducted on the occasion of Aztec emperor Moctezuma II's accession to the throne in 1502. The Franciscan friar Bernardino de Sahagún wrote of witnessing mushroom use in the Florentine Codex (published 1545–1590), and described how some merchants would celebrate upon returning from a successful business trip by consuming mushrooms to evoke revelatory visions. After the defeat of the Aztecs, the Spanish forbade traditional religious practices and rituals that they considered "pagan idolatry", including ceremonial mushroom use. For the next four centuries, the Indians of Mesoamerica hid their use of entheogens from the Spanish authorities. Dozens of species of psychedelic mushrooms are found in Europe, but there is little documented usage of them in Old World history besides the use of Amanita muscaria among Siberian peoples. The few existing accounts that mention psilocybin mushrooms typically lack sufficient information to allow species identification, focusing on their effects. For example, Flemish botanist Carolus Clusius (1526–1609) described the bolond gomba ("crazy mushroom"), used in rural Hungary to prepare love potions. English botanist John Parkinson included details about a "foolish mushroom" in his 1640 herbal Theatricum Botanicum. The first reliably documented report of intoxication with Psilocybe semilanceata—Europe's most common and widespread psychedelic mushroom—involved a British family in 1799, who prepared a meal with mushrooms they had picked in London's Green Park. Modern American banker and amateur ethnomycologist R. Gordon Wasson and his wife, Valentina P. Wasson, a physician, studied the ritual use of psychoactive mushrooms by the native population in the Mazatec village Huautla de Jiménez, Mexico. In 1957, Wasson described the psychedelic visions he experienced during these rituals in "Seeking the Magic Mushroom", an article published in the American weekly Life magazine. Later the same year they were accompanied on a follow-up expedition by French mycologist Roger Heim, who identified several of the mushrooms as Psilocybe species. Heim cultivated the mushrooms in France and sent samples for analysis to Albert Hofmann, a chemist employed by the Swiss pharmaceutical company Sandoz (now Novartis). Hofmann—who had synthesized lysergic acid diethylamide (LSD) in 1938—led a research group that isolated and identified the psychoactive alkaloids psilocybin and psilocin from Psilocybe mexicana, publishing their results in 1958. The team was aided in the discovery process by Hofmann's willingness to ingest mushroom extracts to help verify the presence of the active compounds. Next, Hofmann's team synthesized several structural analogs of these compounds to examine how these structural changes affect psychoactivity. This research led to the development of ethocybin and CZ-74. Because these compounds' physiological effects last only about three and a half hours (about half as long as psilocybin's), they proved more manageable for use in psycholytic therapy. Sandoz also marketed and sold pure psilocybin under the name Indocybin to clinicians and researchers worldwide. There were no reports of serious complications when psilocybin was used in this way. In the early 1960s, Harvard University became a testing ground for psilocybin through the efforts of Timothy Leary and his associates Ralph Metzner and Richard Alpert (who later changed his name to Ram Dass). Leary obtained synthesized psilocybin from Hofmann through Sandoz Pharmaceuticals. Some studies, such as the Concord Prison Experiment, suggested promising results using psilocybin in clinical psychiatry. But according to a 2008 review of safety guidelines in human hallucinogenic research, Leary's and Alpert's well-publicized termination from Harvard and later advocacy of hallucinogen use "further undermined an objective scientific approach to studying these compounds". In response to concerns about the increase in unauthorized use of psychedelic drugs by the general public, psilocybin and other hallucinogenic drugs were unfavorably covered in the press and faced increasingly restrictive laws. In the U.S., laws passed in 1966 that prohibited the production, trade, or ingestion of hallucinogenic drugs; Sandoz stopped producing LSD and psilocybin the same year. In 1970, Congress passed "The Federal Comprehensive Drug Abuse Prevention and Control Act" that made LSD, peyote, psilocybin, and other hallucinogens illegal to use for any purpose, including scientific research. United States politicians' agenda against LSD usage had swept psilocybin along with it into the Schedule I category of illicit drugs. Such restrictions on the use of these drugs in human research made funding for such projects difficult to obtain, and scientists who worked with psychedelic drugs faced being "professionally marginalized". Although Hofmann tested these compounds on himself, he never advocated their legalization or medical use. In his 1979 book LSD—mein Sorgenkind (LSD—My Problem Child), he described the problematic use of these hallucinogens as inebriants. Despite the legal restrictions on psilocybin use, the 1970s witnessed the emergence of psilocybin as the "entheogen of choice". This was due in large part to wide dissemination of information on the topic, which included works such as those by Carlos Castaneda and several books that taught the technique of growing psilocybin mushrooms. One of the most popular of the latter group, Psilocybin: Magic Mushroom Grower's Guide, was published in 1976 under the pseudonyms O. T. Oss and O. N. Oeric by Jeremy Bigwood, Dennis J. McKenna, K. Harrison McKenna, and Terence McKenna. Over 100,000 copies were sold by 1981. As ethnobiologist Jonathan Ott explains, "These authors adapted San Antonio's technique (for producing edible mushrooms by casing mycelial cultures on a rye grain substrate; San Antonio 1971) to the production of Psilocybe [Stropharia] cubensis. The new technique involved the use of ordinary kitchen implements, and for the first time the layperson was able to produce a potent entheogen in his own home, without access to sophisticated technology, equipment or chemical supplies." San Antonio's technique describes a method to grow the common edible mushroom Agaricus bisporus. Because of lack of clarity about laws concerning psilocybin mushrooms, specifically in the form of sclerotia (also known as "truffles"), in the late 1990s and early 2000s European retailers commercialized and marketed them in smartshops in the Netherlands, the UK, and online. Several websites emerged that contributed to the accessibility of information on the mushrooms' description, use, and effects, and users exchanged mushroom experiences. Since 2001, six EU countries have tightened their legislation on psilocybin mushrooms in response to concerns about their prevalence and increasing usage. In the 1990s, hallucinogens and their effects on human consciousness were again the subject of scientific study, particularly in Europe. Advances in neuropharmacology and neuropsychology and the availability of brain imaging techniques have provided impetus for using drugs like psilocybin to probe the "neural underpinnings of psychotic symptom formation including ego disorders and hallucinations". Recent studies in the U.S. have attracted attention from the popular press and brought psilocybin back into the limelight. Reported effects The effects of psilocybin are highly variable and depend on the mindset and environment in which the user has the experience: factors commonly referred to as set and setting. In the early 1960s, Timothy Leary and colleagues at Harvard University investigated the role of set and setting on the effects of psilocybin. They administered the drug to 175 volunteers (from various backgrounds) in an environment intended to be similar to a comfortable living room. 98 of the subjects were given questionnaires to assess their experiences and the contribution of background and situational factors. Individuals who had experience with psilocybin prior to the study reported more pleasant experiences than those for whom the drug was novel. Group size, dosage, preparation, and expectancy were important determinants of the drug response. In general, those in groups of more than eight felt that the groups were less supportive, and their experiences less pleasant. Conversely, smaller groups (fewer than six) were seen as more supportive and reported more positive reactions to the drug in those groups. Leary and colleagues proposed that psilocybin heightens suggestibility, making an individual more receptive to interpersonal interactions and environmental stimuli. These findings were affirmed in a later review by Jos ten Berge (1999), who concluded that dosage, set, and setting were fundamental factors in determining the outcome of experiments that tested the effects of psychedelic drugs on artists' creativity. After ingesting psilocybin, the user may experience a wide range of emotional effects which can include: feelings of disorientation, lethargy, giddiness, euphoria, joy, and depression. In one study, 31% of volunteers given a high dose reported feelings of significant fear and 17% experienced transient paranoia. In studies at Johns Hopkins, among those given a moderate dose (but still enough to "give a high probability of a profound and beneficial experience"), negative experiences were rare, whereas one-third of those given a high dose experienced anxiety or paranoia. Low doses can induce hallucinatory effects. Closed-eye hallucinations may occur, where the affected individual sees multicolored geometric shapes and vivid imaginative sequences. Some individuals report synesthesia, such as tactile sensations when viewing colors. At higher doses, psilocybin can lead to "intensification of affective responses, enhanced ability for introspection, regression to primitive and childlike thinking, and activation of vivid memory traces with pronounced emotional undertones". Open-eye visual hallucinations are common, and may be very detailed, although rarely confused with reality. Physical effects Common responses include pupil dilation (93%); changes in heart rate (100%), including increases (56%), decreases (13%), and variable responses (31%); changes in blood pressure (84%), including hypotension (34%), hypertension (28%), and general instability (22%); changes in stretch reflex (86%), including increases (80%) and decreases (6%); nausea (44%); tremor (25%); and dysmetria (16%) (inability to properly direct or limit motions). The temporary increases in blood pressure caused by the drug can be a risk factor for users with pre-existing hypertension. These qualitative somatic effects caused by psilocybin have been corroborated by several early clinical studies. A 2005 magazine survey of clubgoers in the UK found that nausea or vomiting was experienced by over a quarter of those who had used psilocybin mushrooms in the last year, although this effect is caused by the mushroom rather than psilocybin itself. In one study, administration of gradually increasing dosages of psilocybin daily for 21 days had no measurable effect on electrolyte levels, blood sugar levels, or liver toxicity tests. Psychiatric effects and perceptual disturbances Psilocybin is known to strongly influence the subjective experience of the passage of time. Users often feel as if time is slowed down, resulting in the perception that "minutes appear to be hours" or "time is standing still". Studies have demonstrated that psilocybin significantly impairs subjects' ability to gauge time intervals longer than 2.5 seconds, impairs their ability to synchronize to inter-beat intervals longer than 2 seconds, and reduces their preferred tapping rate. These results are consistent with the drug's role in affecting prefrontal cortex activity, and the role that the prefrontal cortex is known to play in time perception. However, the neurochemical basis of psilocybin's effects on the perception of time are not known with certainty. Users having a pleasant experience can feel a sense of connection to others, nature, and the universe; other perceptions and emotions are also often intensified. Users having an unpleasant experience (a "bad trip") describe a reaction accompanied by fear, other unpleasant feelings, and occasionally by dangerous behavior. In general, the phrase "bad trip" is used to describe a reaction that is characterized primarily by fear or other unpleasant emotions, not just a transitory experience of such feelings. A variety of factors may contribute to a psilocybin user experiencing a bad trip, including "tripping" during an emotional or physical low or in a non-supportive environment (see: set and setting). Ingesting psilocybin in combination with other drugs, including alcohol, can also increase the likelihood of a bad trip. Other than the duration of the experience, the effects of psilocybin are similar to comparable dosages of lysergic acid diethylamide (LSD) or mescaline. However, in the Psychedelics Encyclopedia, author Peter Stafford noted, "The psilocybin experience seems to be warmer, not as forceful and less isolating. It tends to build connections between people, who are generally much more in communication than when they use LSD." Theory of mind network and default mode network Psychedelics, including psilocybin, have been shown to affect different clusters of brain regions known as the "theory of mind network" (ToMN) and the default mode network (DMN). The ToMN involves making inferences and understanding social situations based on patterns whereas, the DMN relates more to introspection and one's sense of self. The DMN in particular is related to increased rumination and worsening self-image in patients with major depressive disorder (MDD). In studies done with single use psilocybin, areas of the DMN showed decreased functional connectivity (communication between areas of the brain). This provides functional insight into the work of psilocybin in increasing one's sense of connection to one's surroundings, as the areas of the brain involved in introspection decrease in functionality under the effects of the drug. Conversely, areas of the brain involved in the ToMN showed increased activity and functional activation in response to psychedelics. These results were not unique to psilocybin and there was no significant difference in brain activation found in similar trials of mescaline and LSD. Information and studies into the DMN and ToMN are relatively sparse and their connections to other psychiatric illnesses and the use of psychedelics is still largely unknown. Group perceptions Through further anthropological studies regarding "personal insights" and the psychosocial effects of psilocybin, it can be seen in many traditional societies that powerful mind-active substances such as psilocybin are regularly "consumed ritually for therapeutic purposes or for transcending normal, everyday reality". Positive effects that psilocybin has on individuals can be observed by taking on an anthropological approach and moving away from the Western biomedical view; this is aided by the studies done by Leary. Within certain traditional societies, where the use of psilocybin is frequent for shamanic healing rituals, group collectives praise their guide, healer and shaman for helping alleviate their pains, aches and hurt. They do this through a group ritual practice where the group, or just the guide, ingests psilocybin to help extract any "toxic psychic residues or sorcerous implants" found in one's body. Group therapies using "classic" psychedelics are becoming more commonly used in the Western world in clinical practice. This is speculated to grow, provided the evidence remains indicative of their safety and efficacy. In social sense, the group is shaped by their experiences surrounding psilocybin and how they view the plant collectively. As mentioned in the anthropology article, the group partakes in a "journey" together, thus adding to the spiritual, social body where roles, hierarchies and gender are subjectively understood. Adverse effects Most of the comparatively few fatal incidents that are associated with psychedelic mushroom usage involve the simultaneous use of other drugs, especially alcohol. A common adverse effect resulting from psilocybin mushroom use involves "bad trips" or panic reactions, in which affected individuals become anxious, confused, agitated, or disoriented. Accidents, self-injury, or suicide attempts can result from serious cases of acute psychotic episodes. Although no studies have linked psilocybin with birth defects, it is recommended that pregnant women avoid its usage. Toxicity Psilocybin has low toxicity, indicating that it has low potential for inducing life-threatening events like breathing or heart problems. Research shows that health risks may develop with use of psilocybin. Nonetheless, hospitalizations from it are rare, and overdoses are generally mild and self-limiting. A review regarding the management of psychedelic overdoses suggested that psilocybin-related overdose management should prioritize managing the immediate adverse effects, such as anxiety and paranoia, rather than specific pharmacological interventions, as psilocybin's physiological toxicity tends to be rather limited. One analysis of people hospitalized from psilocybin poisoning found high urine concentrations of phenethylamine (PEA), indicating that PEA may contribute to the effects of psilocybin poisoning. In rats, the median lethal dose (LD50) when administered orally is 280 milligrams per kilogram (mg/kg), approximately one and a half times that of caffeine. The lethal dose of psilocybin when administered intravenously in mice is 285mg/kg and in rats is 280 mg/kg. When administered intravenously in rabbits, psilocybin's LD50 is approximately 12.5 mg/kg. Psilocybin comprises approximately 1% of the weight of Psilocybe cubensis mushrooms, and so nearly of dried mushrooms, or of fresh mushrooms, would be required for a person to reach the 280 mg/kg LD50 value of rats. Based on the results of animal studies, the lethal dose of psilocybin has been extrapolated to be 6 grams, 1000 times greater than the effective dose of 6 milligrams. The Registry of Toxic Effects of Chemical Substances assigns psilocybin a relatively high therapeutic index of 641 (higher values correspond to a better safety profile); for comparison, the therapeutic indices of aspirin and nicotine are 199 and 21, respectively. The lethal dose from psilocybin toxicity alone is unknown, and has rarely been documented—, only two cases attributed to overdosing on hallucinogenic mushrooms (without concurrent use of other drugs) have been reported in the scientific literature and may involve other factors aside from psilocybin. Psychiatric Panic reactions can occur after consumption of psilocybin-containing mushrooms, especially if the ingestion is accidental or otherwise unexpected. Reactions characterized by violent behavior, suicidal thoughts, schizophrenia-like psychosis, and convulsions have been reported in the literature. A 2005 survey, conducted in the United Kingdom, found that almost a quarter of those who had used psilocybin mushrooms in the past year had experienced a panic attack. Other adverse effects, less frequently reported, include paranoia, confusion, prolonged derealization (disconnection from reality), and mania. Psilocybin usage can temporarily induce a state of depersonalization disorder. Usage by those with schizophrenia can induce acute psychotic states requiring hospitalization. The similarity of psilocybin-induced symptoms to those of schizophrenia has made the drug a useful research tool in behavioral and neuroimaging studies of this psychotic disorder. In both cases, psychotic symptoms are thought to arise from a "deficient gating of sensory and cognitive information" in the brain that ultimately lead to "cognitive fragmentation and psychosis". Flashbacks (spontaneous recurrences of a previous psilocybin experience) can occur long after having used psilocybin mushrooms. Hallucinogen persisting perception disorder (HPPD) is characterized by a continual presence of visual disturbances similar to those generated by psychedelic substances. Neither flashbacks nor HPPD are commonly associated with psilocybin usage, and correlations between HPPD and psychedelics are further obscured by polydrug use and other variables. Tolerance and dependence Tolerance to psilocybin builds and dissipates quickly; ingesting psilocybin more than about once a week can lead to diminished effects. Tolerance dissipates after a few days, so doses can be spaced several days apart to avoid the effect. A cross-tolerance can develop between psilocybin and the pharmacologically similar LSD, and between psilocybin and phenethylamines such as mescaline and DOM. Repeated use of psilocybin does not lead to physical dependence. A 2008 study concluded that, based on US data from the period 2000–2002, adolescent-onset (defined here as ages 11–17) usage of hallucinogenic drugs (including psilocybin) did not increase the risk of drug dependence in adulthood; this was in contrast to adolescent usage of cannabis, cocaine, inhalants, anxiolytic medicines, and stimulants, all of which were associated with "an excess risk of developing clinical features associated with drug dependence". Likewise, a 2010 Dutch study ranked the relative harm of psilocybin mushrooms compared to a selection of 19 recreational drugs, including alcohol, cannabis, cocaine, ecstasy, heroin, and tobacco. Psilocybin mushrooms were ranked as the illicit drug with the lowest harm, corroborating conclusions reached earlier by expert groups in the United Kingdom. Cultural significance and "mystical" experiences Psilocybin mushrooms have been and continue to be used in Indigenous American cultures in religious, divinatory, or spiritual contexts. Reflecting the meaning of the word entheogen ("the god within"), the mushrooms are revered as powerful spiritual sacraments that provide access to sacred worlds. Typically used in small group community settings, they enhance group cohesion and reaffirm traditional values. Terence McKenna documented the worldwide practices of psilocybin mushroom usage as part of a cultural ethos relating to the Earth and mysteries of nature, and suggested that mushrooms enhanced self-awareness and a sense of contact with a "Transcendent Other"—reflecting a deeper understanding of our connectedness with nature. Psychedelic drugs can induce states of consciousness that have lasting personal meaning and spiritual significance in individuals who are religious or spiritually inclined; these states are called mystical experiences. Some scholars have proposed that many of the qualities of a drug-induced mystical experience are indistinguishable from mystical experiences achieved through non-drug techniques, such as meditation or holotropic breathwork. In the 1960s, Walter Pahnke and colleagues systematically evaluated mystical experiences (which they called "mystical consciousness") by categorizing their common features. These categories, according to Pahnke, "describe the core of a universal psychological experience, free from culturally determined philosophical or theological interpretations", and allow researchers to assess mystical experiences on a qualitative, numerical scale. In the 1962 Marsh Chapel Experiment, which was run by Pahnke at the Harvard Divinity School under the supervision of Timothy Leary, almost all of the graduate degree divinity student volunteers who received psilocybin reported profound religious experiences. One of the participants was religious scholar Huston Smith, author of several textbooks on comparative religion; he later described his experience as "the most powerful cosmic homecoming I have ever experienced." In a 25-year followup to the experiment, all of the subjects given psilocybin described their experience as having elements of "a genuine mystical nature and characterized it as one of the high points of their spiritual life". Psychedelic researcher Rick Doblin considered the study partially flawed due to incorrect implementation of the double-blind procedure, and several imprecise questions in the mystical experience questionnaire. Nevertheless, he said that the study cast "a considerable doubt on the assertion that mystical experiences catalyzed by drugs are in any way inferior to non-drug mystical experiences in both their immediate content and long-term effects". This sentiment was echoed by psychiatrist William A. Richards, who in a 2007 review stated "[psychedelic] mushroom use may constitute one technology for evoking revelatory experiences that are similar, if not identical, to those that occur through so-called spontaneous alterations of brain chemistry." A group of researchers from Johns Hopkins School of Medicine led by Roland Griffiths conducted a study to assess the immediate and long-term psychological effects of the psilocybin experience, using a modified version of the mystical experience questionnaire and a rigorous double-blind procedure. When asked in an interview about the similarity of his work with Leary's, Griffiths explained the difference: "We are conducting rigorous, systematic research with psilocybin under carefully monitored conditions, a route which Dr. Leary abandoned in the early 1960s." The National Institute of Drug Abuse-funded study, published in 2006, has been praised by experts for the soundness of its experimental design. In the experiment, 36 volunteers without prior experience with hallucinogens were given psilocybin and methylphenidate (Ritalin) in separate sessions; the methylphenidate sessions served as a control and psychoactive placebo. The degree of mystical experience was measured using a questionnaire developed by Ralph W. Hood; 61% of subjects reported a "complete mystical experience" after their psilocybin session, while only 13% reported such an outcome after their experience with methylphenidate. Two months after taking psilocybin, 79% of the participants reported moderately to greatly increased life satisfaction and sense of well-being. About 36% of participants also had a strong to extreme "experience of fear" or dysphoria (i.e., a "bad trip") at some point during the psilocybin session (which was not reported by any subject during the methylphenidate session); about one-third of these (13% of the total) reported that this dysphoria dominated the entire session. These negative effects were reported to be easily managed by the researchers and did not have a lasting negative effect on the subject's sense of well-being. A follow-up study conducted 14 months after the original psilocybin session confirmed that participants continued to attribute deep personal meaning to the experience. Almost one-third of the subjects reported that the experience was the single most meaningful or spiritually significant event of their lives, and over two-thirds reported it among their five most spiritually significant events. About two-thirds indicated that the experience increased their sense of well-being or life satisfaction. Even after 14 months, those who reported mystical experiences scored on average 4 percentage points higher on the personality trait of Openness/Intellect; personality traits are normally stable across the lifespan for adults. Likewise, in a recent (2010) web-based questionnaire study designed to investigate user perceptions of the benefits and harms of hallucinogenic drug use, 60% of the 503 psilocybin users reported that their use of psilocybin had a long-term positive impact on their sense of well-being. While many recent studies have concluded that psilocybin can cause mystical-type experiences having substantial and sustained personal meaning and spiritual significance, not all the medical community agree. Paul R. McHugh, formerly director of the Department of Psychiatry and Behavioral Science at Johns Hopkins, responded as follows in a book review: "The unmentioned fact in The Harvard Psychedelic Club is that LSD, psilocybin, mescaline, and the like produce not a "higher consciousness" but rather a particular kind of "lower consciousness" known well to psychiatrists and neurologists—namely, 'toxic delirium.'" Available forms Although psilocybin may be prepared synthetically, outside of the research setting it is not typically used in this form. The psilocybin present in certain species of mushrooms can be ingested in several ways: by consuming fresh or dried fruit bodies, by preparing an herbal tea, or by combining with other foods to mask the bitter taste. In rare cases people have injected mushroom extracts intravenously. Chemistry Physical properties Psilocybin is a naturally-occurring substituted tryptamine that features an indole ring linked to an aminoethyl substituent. It is structurally related to serotonin, a monoamine neurotransmitter which is a derivative of the amino acid tryptophan. Psilocybin is a member of the general class of tryptophan-based compounds that originally functioned as antioxidants in earlier life forms before assuming more complex functions in multicellular organisms, including humans. Other related indole-containing psychedelic compounds include dimethyltryptamine, found in many plant species and in trace amounts in some mammals, and bufotenin, found in the skin of certain amphibians, especially the Colorado River toad. Psilocybin is a white, crystalline solid that is soluble in water, methanol and ethanol but insoluble in nonpolar organic solvents such as chloroform and petroleum ether. It has a melting point between , and an ammonia-like taste. Its pKa values are estimated to be 1.3 and 6.5 for the two successive phosphate hydroxy groups and 10.4 for the dimethylamine nitrogen, so it typically exists as a zwitterionic structure. There are two known crystalline polymorphs of psilocybin, as well as reported hydrated phases. Psilocybin rapidly oxidizes upon exposure to light—an important consideration when using it as an analytical standard. Laboratory synthesis Albert Hofmann et al. were the first team to synthesize psilocybin in 1958. Since that time, various chemists have improved the methods for the laboratory synthesis and purification of psilocybin. In particular, Shirota et al. reported a novel method in 2003 for the synthesis of psilocybin at the gram scale from 4-hydroxyindole that does not require chromatographic purification. Fricke et al. described an enzymatic pathway for the synthesis of psilocybin and psilocin, publishing their results in 2017. Sherwood et al. significantly improved upon Shirota's method (producing at the kilogram scale while employing less expensive reagents), publishing their results in 2020. Biosynthesis Isotopic labeling experiments from the 1960s suggested that the biosynthesis of psilocybin was a four-step process: decarboxylation of tryptophan to tryptamine N,N-dimethylation of tryptamine at the N9 position to dimethyltryptamine 4-hydroxylation of dimethyltryptamine to psilocin O-phosphorylation of psilocin to psilocybin More recent research has demonstrated that—at least in P. cubensis—O-phosphorylation is in fact the third step, and that neither dimethyltryptamine nor psilocin are intermediates. The sequence of the intermediate steps has been shown to involve four enzymes (PsiD, PsiH, PsiK, and PsiM) in P. cubensis and P. cyanescens, although it is possible that the biosynthetic pathway differs between species. These enzymes are encoded in gene clusters in Psilocybe, Panaeolus, and Gymnopilus. Escherichia coli has been genetically modified to manufacture large amounts of psilocybin. Psilocybin can be produced de novo in GM yeast. Pharmacology Pharmacodynamics Psilocybin is a psychoplastogen, which refers to a compound capable of promoting rapid and sustained neuroplasticity. Psilocybin is rapidly dephosphorylated in the body to psilocin, which is an agonist for several serotonin receptors, which are also known as 5-hydroxytryptamine (5-HT) receptors. In rats, psilocin binds with high affinity to 5-HT2A receptors and low affinity to 5-HT1 receptors, including 5-HT1A and 5-HT1D; effects are also mediated via 5-HT2C receptors. The psychotomimetic (mimicking the mind distortion present in psychosis) effects of psilocin can be blocked in a dose-dependent fashion by the 5-HT2A antagonist drug ketanserin. Various lines of evidence have shown that interactions with non-5-HT2 receptors also contribute to the subjective and behavioral effects of the drug. For example, psilocin indirectly increases the concentration of the neurotransmitter dopamine in the basal ganglia, and some psychotomimetic symptoms of psilocin are reduced by haloperidol, a non-selective dopamine receptor antagonist. Taken together, these suggest that there may be an indirect dopaminergic contribution to psilocin's psychotomimetic effects. Psilocybin and psilocin have no affinity for dopamine D2 receptor, unlike another common 5-HT receptor agonist, lysergic acid diethylamide (LSD). Psilocin antagonizes histamine H1 receptors with moderate affinity, compared to LSD which has lower affinity. Serotonin receptors are located in numerous parts of the brain, including the cerebral cortex, and are involved in a wide range of functions, including regulation of mood, motivation, body temperature, appetite and libido. Psilocybin induces region-dependent alterations in glutamate that may be associated with subjective experiences of ego dissolution. Pharmacokinetics The effects of the drug begin 10–40 minutes after ingestion, and last 2–6 hours depending on dose, species, and individual metabolism. The half life of psilocybin is 163 ± 64 minutes when taken orally, or 74.1 ± 19.6 minutes when injected intravenously. Psilocybin is metabolized mostly in the liver. As it becomes converted to psilocin, it undergoes a first-pass effect, whereby its concentration is greatly reduced before it reaches the systemic circulation. Psilocin is broken down by the enzyme monoamine oxidase (MAO) to produce several metabolites that can circulate in the blood plasma, including 4-hydroxyindole-3-acetaldehyde (4-HIAL), 4-hydroxyindole-3-acetic acid (4-HIAA), and 4-hydroxytryptophol (4-HTOL). Some psilocin is not broken down by enzymes and instead forms a glucuronide; this is a biochemical mechanism animals use to eliminate substances by linking them with glucuronic acid, which can then be excreted in the urine. Psilocin is glucuronidated by the glucuronosyltransferase enzymes UGT1A9 in the liver, and by UGT1A10 in the small intestine. Based on studies using animals, about 50% of ingested psilocybin is absorbed through the stomach and intestine. About 80% of psilocin is glucuronidated into psilocin-O-glucuronide and about 4% is demethylated and oxidatively deaminated via MAO into 4-HIAL. 4-HIAL is subsequently oxidized by aldehyde dehydrogenase (ALDH) into 4-HIAA or can be converted by alcohol dehydrogenase (ADH) into 4-HTOL. Within 24 hours, about 65% of the absorbed psilocybin is excreted into the urine, and a further 15–20% is excreted in the bile and feces. Although most of the remaining drug is eliminated in this way within 8 hours, it is still detectable in the urine after 7 days. Clinical studies show that psilocin concentrations in the plasma of adults average about 8 μg/liter within 2 hours after ingestion of a single 15 mg oral psilocybin dose; psychological effects occur with a blood plasma concentration of 4–6 μg/liter. Psilocybin is approximately 1/100 the potency of LSD on a weight per weight basis, and the physiological effects last about half as long. Monoamine oxidase inhibitors (MAOI) have been known to prolong and enhance the effects of dimethyltryptamine (DMT) and one study assumed that the effect on psilocybin would be similar since it is a structural analogue of DMT. However, only a small portion of psilocin appears to be metabolized by MAO. Alcohol consumption may enhance the effects of psilocybin, because acetaldehyde, one of the primary breakdown metabolites of consumed alcohol, reacts with biogenic amines present in the body to produce MAOIs related to tetrahydroisoquinoline and β-carboline. Tobacco smokers may also experience more powerful effects with psilocybin, because tobacco smoke exposure decreases the activity of MAO in the brain and peripheral organs. Analytical methods Several relatively simple chemical tests—commercially available as reagent testing kits—can be used to assess the presence of psilocybin in extracts prepared from mushrooms. The drug reacts in the Marquis test to produce a yellow color, and a green color in the Mandelin reagent. Neither of these tests, however, is specific for psilocybin; for example, the Marquis test will react with many classes of controlled drugs, such as those containing primary amino groups and unsubstituted benzene rings, including amphetamine and methamphetamine. Ehrlich's reagent and DMACA reagent are used as chemical sprays to detect the drug after thin layer chromatography. Many modern techniques of analytical chemistry have been used to quantify psilocybin levels in mushroom samples. Although the earliest methods commonly used gas chromatography, the high temperature required to vaporize the psilocybin sample prior to analysis causes it to spontaneously lose its phosphoryl group and become psilocin—making it difficult to chemically discriminate between the two drugs. In forensic toxicology, techniques involving gas chromatography coupled to mass spectrometry (GC–MS) are the most widely used due to their high sensitivity and ability to separate compounds in complex biological mixtures. These techniques include ion mobility spectrometry, capillary zone electrophoresis, ultraviolet spectroscopy, and infrared spectroscopy. High-performance liquid chromatography (HPLC) is used with ultraviolet, fluorescence, electrochemical, and electrospray mass spectrometric detection methods. Various chromatographic methods have been developed to detect psilocin in body fluids: the rapid emergency drug identification system (REMEDi HS), a drug screening method based on HPLC; HPLC with electrochemical detection; GC–MS; and liquid chromatography coupled to mass spectrometry. Although the determination of psilocin levels in urine can be performed without sample clean-up (i.e., removing potential contaminants that make it difficult to accurately assess concentration), the analysis in plasma or serum requires a preliminary extraction, followed by derivatization of the extracts in the case of GC–MS. A specific immunoassay has also been developed to detect psilocin in whole blood samples. A 2009 publication reported using HPLC to quickly separate forensically important illicit drugs including psilocybin and psilocin, which were identifiable within about half a minute of analysis time. These analytical techniques to determine psilocybin concentrations in body fluids are, however, not routinely available, and not typically used in clinical settings. Natural occurrence Psilocybin is present in varying concentrations in over 200 species of Basidiomycota mushrooms. In a 2000 review on the worldwide distribution of hallucinogenic mushrooms, Gastón Guzmán and colleagues considered these to be distributed amongst the following genera: Psilocybe (116 species), Gymnopilus (14), Panaeolus (13), Copelandia (12), Hypholoma (6), Pluteus (6), Inocybe (6), Conocybe (4), Panaeolina (4), Gerronema (2), and Galerina (1 species). Guzmán increased his estimate of the number of psilocybin-containing Psilocybe to 144 species in a 2005 review. The majority of these are found in Mexico (53 species), with the remainder distributed in the United States and Canada (22), Europe (16), Asia (15), Africa (4), and Australia and associated islands (19). The diversity of psilocybin mushrooms is reported to have been increased by horizontal transfer of the psilocybin gene cluster between unrelated mushroom species. In general, psilocybin-containing species are dark-spored, gilled mushrooms that grow in meadows and woods of the subtropics and tropics, usually in soils rich in humus and plant debris. Psilocybin mushrooms occur on all continents, but the majority of species are found in subtropical humid forests. Psilocybe species commonly found in the tropics include P. cubensis and P. subcubensis. P. semilanceata—considered by Guzmán to be the world's most widely distributed psilocybin mushroom—is found in Europe, North America, Asia, South America, Australia and New Zealand, but is entirely absent from Mexico. Although the presence or absence of psilocybin is not of much use as a chemotaxonomical marker at the familial level or higher, it is used to classify taxa of lower taxonomic groups. Both the caps and the stems contain psychoactive compounds, although the caps consistently contain more. The spores of these mushrooms do not contain psilocybin or psilocin. The total potency varies greatly between species and even between specimens of a species collected or grown from the same strain. Because most psilocybin biosynthesis occurs early in the formation of fruit bodies or sclerotia, younger, smaller mushrooms tend to have a higher concentration of the drug than larger, mature mushrooms. In general, the psilocybin content of mushrooms is quite variable (ranging from almost nothing to 2.5% of the dry weight) and depends on species, strain, growth and drying conditions, and mushroom size. Cultivated mushrooms have less variability in psilocybin content than wild mushrooms. The drug is more stable in dried than fresh mushrooms; dried mushrooms retain their potency for months or even years, while mushrooms stored fresh for four weeks contain only traces of the original psilocybin. The psilocybin contents of dried herbarium specimens of Psilocybe semilanceata in one study were shown to decrease with the increasing age of the sample: collections dated 11, 33, or 118 years old contained 0.84%, 0.67%, and 0.014% (all dry weight), respectively. Mature mycelia contain some psilocybin, while young mycelia (recently germinated from spores) lack appreciable amounts. Many species of mushrooms containing psilocybin also contain lesser amounts of the analog compounds baeocystin and norbaeocystin, chemicals thought to be biogenic precursors. Although most species of psilocybin-containing mushrooms bruise blue when handled or damaged due to the oxidization of phenolic compounds, this reaction is not a definitive method of identification or determining a mushroom's potency. Societal perception and current usage Legal status Advocacy for tolerance Despite being illegal in many typically Western countries, such as the UK, Australia and some US states, less conservative governments opt to nurture the legal use of psilocybin and other psychedelic drugs. In Amsterdam, Netherlands, authorities provide education and promotion on the safe use of psychedelic drugs, such as psilocybin, in an aim to reduce public harm. Similarly, religious groups like America's Uniao do Vegetal, UDV, use psychedelics in traditional ceremonies. A report from the U.S. Government Accountability Office (GAO) notes that people may petition the DEA for exemptions to use psilocybin for religious purposes. From 1 July 2023, the Australian medicines regulator has permitted psychiatrists to prescribe psilocybin for the therapeutic treatment of treatment-resistant depression. Advocates for legalization argue there is a lack of evidence of harm, and potential use in treating certain mental health conditions. Research is difficult to conduct because of the legal status of psychoactive substances. Advocates for legalization also promote the utility of "ego dissolution" and argue bans are cultural discrimination against traditional users. Usage A 2009 national survey of drug use by the US Department of Health and Human Services concluded that the number of first-time psilocybin mushroom users in the United States was roughly equivalent to the number of first-time users of cannabis. A June 2024 report by the RAND Corporation suggests the total number of use days for psychedelics is two orders of magnitude smaller than it is for cannabis, and unlike people who use cannabis and many other drugs, infrequent users of psychedelics account for most of the total days of use. The 2024 report by the RAND Corporation suggests psilocybin mushrooms may be the most prevalent psychedelic drug among adults in the United States. In European countries, the lifetime prevalence estimates of psychedelic mushroom usage among young adults (15–34 years) range from 0.3% to 14.1%. In modern Mexico, traditional ceremonial use survives among several indigenous groups, including the Nahuas, the Matlatzinca, the Totonacs, the Mazatecs, Mixes, Zapotecs, and the Chatino. Although hallucinogenic Psilocybe species are abundant in low-lying areas of Mexico, most ceremonial use takes places in mountainous areas of elevations greater than . Guzmán suggests this is a vestige of Spanish colonial influence from several hundred years earlier, when mushroom use was persecuted by the Catholic Church. Research Psilocybin has been a subject of clinical research since the early 1960s, when the Harvard Psilocybin Project evaluated the potential value of psilocybin as a treatment for certain personality disorders. Beginning in the 2000s, psilocybin has been investigated for its possible role in the treatment of nicotine dependence, alcohol dependence, obsessive–compulsive disorder (OCD), cluster headache, cancer-related existential distress, anxiety disorders, and certain mood disorders. In 2018, the United States Food and Drug Administration (FDA) granted breakthrough therapy designation for psilocybin-assisted therapy for treatment-resistant depression. A systematic review published in 2021 found that the use of psilocybin as a pharmaceutical substance was associated with reduced intensity of depression symptoms. The role of psilocybin as a possible psychoplastogen is also being examined. Depression Clinical trials, including both open-label trials and double-blind randomized controlled trials, have found that single doses of psilocybin produce rapid and long-lasting antidepressant effects outperforming placebo in people with major depressive disorder and treatment-resistant depression. Combined with brief psychological support in a phase 2 trial, it has been found to produce dose-dependent improvements in depressive symptoms, with 25mg (a moderate dose) more effective than 10mg (a low dose) and 10mg more effective than 1mg (non-psychoactive and equivalent to placebo). The antidepressant effects of psilocybin with psychological support have been found to last at least 6weeks following a single dose. However, some trials have not found psilocybin to significantly outperform placebo in the treatment of depression. In addition, a phase 2 trial found that two 25mg doses of psilocybin 3weeks apart versus daily treatment with the selective serotonin reuptake inhibitor (SSRI) escitalopram (Lexapro) for 6weeks (plus two putatively non-psychoactive 1mg doses of psilocybin 3weeks apart) did not show a statistically significant difference in reduction of depressive symptoms between groups. However, reductions in depressive symptoms were numerically greater with psilocybin, some secondary measures favored psilocybin, and the rate of remission was statistically higher with psilocybin (57% with psilocybin vs. 28% with escitalopram). In any case, the antidepressant effect size of psilocybin over escitalopram appears to be small. Functional unblinding by their psychoactive effects and positive psychological expectancy effects (i.e., the placebo effect) are major limitations and sources of bias of clinical trials of psilocybin and other psychedelics for treatment of depression. In addition, as of September 2024, psilocybin and other psychedelics (excluding MDMA) have only been assessed in up to phase 2 clinical trials for psychiatric disorders and have not yet completed larger and more rigorous phase 3 trials or received regulatory approval for medical use. A potential risk of frequent repeated use of psilocybin and other serotonergic psychedelics for psychiatric disorders is cardiac fibrosis and valvulopathy caused by serotonin 5-HT2B receptor activation. However, single high doses or widely spaced doses (e.g., months) are widely thought to be safe and concerns about cardiac toxicity apply more to chronic psychedelic microdosing or very frequent use (e.g., weekly).	Biology and health sciences	Recreational drugs	Health
38481	https://en.wikipedia.org/wiki/Human%20voice	Human voice	The human voice consists of sound made by a human being using the vocal tract, including talking, singing, laughing, crying, screaming, shouting, humming or yelling. The human voice frequency is specifically a part of human sound production in which the vocal folds (vocal cords) are the primary sound source. (Other sound production mechanisms produced from the same general area of the body involve the production of unvoiced consonants, clicks, whistling and whispering.) Generally speaking, the mechanism for generating the human voice can be subdivided into three parts; the lungs, the vocal folds within the larynx (voice box), and the articulators. The lungs, the "pump" must produce adequate airflow and air pressure to vibrate vocal folds. The vocal folds (vocal cords) then vibrate to use airflow from the lungs to create audible pulses that form the laryngeal sound source. The muscles of the larynx adjust the length and tension of the vocal folds to 'fine-tune' pitch and tone. The articulators (the parts of the vocal tract above the larynx consisting of tongue, palate, cheek, lips, etc.) articulate and filter the sound emanating from the larynx and to some degree can interact with the laryngeal airflow to strengthen or weaken it as a sound source. The vocal folds, in combination with the articulators, are capable of producing highly intricate arrays of sound. The tone of voice may be modulated to suggest emotions such as anger, surprise, fear, happiness or sadness. The human voice is used to express emotion, and can also reveal the age and sex of the speaker. Singers use the human voice as an instrument for creating music. Voice types and the folds (cords) themselves Adult men and women typically have different sizes of vocal fold; reflecting the male–female differences in larynx size. Adult male voices are usually lower-pitched and have larger folds. The male vocal folds (which would be measured vertically in the opposite diagram), are between 17 mm and 25 mm in length. The female vocal folds are between 12.5 mm and 17.5 mm in length. The folds are within the larynx. They are attached at the back (side nearest the spinal cord) to the arytenoids cartilages, and at the front (side under the chin) to the thyroid cartilage. They have no outer edge as they blend into the side of the breathing tube (the illustration is out of date and does not show this well) while their inner edges or "margins" are free to vibrate (the hole). They have a three layer construction of an epithelium, vocal ligament, then muscle (vocalis muscle), which can shorten and bulge the folds. They are flat triangular bands and are pearly white in color. Above both sides of the vocal cord is the vestibular fold or false vocal cord, which has a small sac between its two folds. The difference in vocal fold size between men and women means that they have differently pitched voices. There is also genetic variation amongst the same sex, with men's and women's singing voices being categorized into types. For example, among men, there are bass, bass-baritone, baritone, baritenor, tenor and countertenor (ranging from E2 to C♯7 and higher), and among women, contralto, alto, mezzo-soprano and soprano (ranging from F3 to C6 and higher). There are additional categories for operatic voices, see voice type. This is not the only source of difference between male and female voice. Men, generally speaking, have a larger vocal tract, which essentially gives the resultant voice a lower-sounding timbre. This is mostly independent of the vocal folds themselves. Voice modulation in spoken language Human spoken language makes use of the ability of almost all people in a given society to dynamically modulate certain parameters of the laryngeal voice source in a consistent manner. The most important communicative, or phonetic, parameters are the voice pitch (determined by the vibratory frequency of the vocal folds) and the degree of separation of the vocal folds, referred to as vocal fold adduction (coming together) or abduction (separating). The ability to vary the ab/adduction of the vocal folds quickly has a strong genetic component, since vocal fold adduction has a life-preserving function in keeping food from passing into the lungs, in addition to the covering action of the epiglottis. Consequently, the muscles that control this action are among the fastest in the body. Children can learn to use this action consistently during speech at an early age, as they learn to speak the difference between utterances such as "apa" (having an abductory–adductory gesture for the p) as "aba" (having no abductory–adductory gesture). They can learn to do this well before the age of two by listening only to the voices of adults around them who have voices much different from their own, and even though the laryngeal movements causing these phonetic differentiations are deep in the throat and not visible to them. If an abductory movement or adductory movement is strong enough, the vibrations of the vocal folds will stop (or not start). If the gesture is abductory and is part of a speech sound, the sound will be called voiceless. However, voiceless speech sounds are sometimes better identified as containing an abductory gesture, even if the gesture was not strong enough to stop the vocal folds from vibrating. This anomalous feature of voiceless speech sounds is better understood if it is realized that it is the change in the spectral qualities of the voice as abduction proceeds that is the primary acoustic attribute that the listener attends to when identifying a voiceless speech sound, and not simply the presence or absence of voice (periodic energy). An adductory gesture is also identified by the change in voice spectral energy it produces. Thus, a speech sound having an adductory gesture may be referred to as a "glottal stop" even if the vocal fold vibrations do not entirely stop. Other aspects of the voice, such as variations in the regularity of vibration, are also used for communication, and are important for the trained voice user to master, but are more rarely used in the formal phonetic code of a spoken language. Physiology and vocal timbre The sound of each individual's voice is thought to be entirely unique not only because of the actual shape and size of an individual's vocal cords but also due to the size and shape of the rest of that person's body, especially the vocal tract, and the manner in which the speech sounds are habitually formed and articulated. (It is this latter aspect of the sound of the voice that can be mimicked by skilled performers.) Humans have vocal folds that can loosen, tighten, or change their thickness, and over which breath can be transferred at varying pressures. The shape of chest and neck, the position of the tongue, and the tightness of otherwise unrelated muscles can be altered. Any one of these actions results in a change in pitch, volume, timbre, or tone of the sound produced. Sound also resonates within different parts of the body, and an individual's size and bone structure can affect somewhat the sound produced by an individual. Singers can also learn to project sound in certain ways so that it resonates better within their vocal tract. This is known as vocal resonation. Another major influence on vocal sound and production is the function of the larynx, which people can manipulate in different ways to produce different sounds. These different kinds of laryngeal function are described as different kinds of vocal registers. The primary method for singers to accomplish this is through the use of the singer's formant, which has been shown to be a resonance added to the normal resonances of the vocal tract above the frequency range of most instruments and so enables the singer's voice to carry better over musical accompaniment. Vocal registration Vocal registration refers to the system of vocal registers within the human voice. A register in the human voice is a particular series of tones, produced in the same vibratory pattern of the vocal folds, and possessing the same quality. Registers originate in laryngeal functioning. They occur because the vocal folds are capable of producing several different vibratory patterns. Each of these vibratory patterns appears within a particular Vocal range of pitches and produces certain characteristic sounds. The occurrence of registers has also been attributed to effects of the acoustic interaction between the vocal fold oscillation and the vocal tract. The term register can be somewhat confusing as it encompasses several aspects of the human voice. The term register can be used to refer to any of the following: A particular part of the vocal range such as the upper, middle, or lower registers. A resonance area such as chest voice or head voice. A phonatory process. A certain vocal timbre. A region of the voice that is defined or delimited by vocal breaks. A subset of a language used for a particular purpose or in a particular social setting. In linguistics, a register language is a language that combines tone and vowel phonation into a single phonological system. Within speech pathology, the term vocal register has three constituent elements: a certain vibratory pattern of the vocal folds, a certain series of pitches, and a certain type of sound. Speech pathologists identify four vocal registers based on the physiology of laryngeal function: the vocal fry register, the modal register, the falsetto register, and the whistle register. This view is also adopted by many vocal pedagogists. Vocal resonation Vocal resonation is the process by which the basic product of phonation is enhanced in timbre and/or intensity by the air-filled cavities through which it passes on its way to the outside air. Various terms related to the resonation process include amplification, enrichment, enlargement, improvement, intensification, and prolongation; although in strictly scientific usage acoustic authorities would question most of them. The main point to be drawn from these terms by a singer or speaker is that the result of resonation is, or should be, to make a better sound. There are seven areas that may be listed as possible vocal resonators. In sequence from the lowest within the body to the highest, these areas are the chest, the tracheal tree, the larynx itself, the pharynx, the oral cavity, the nasal cavity, and the sinuses. Influences of the human voice The twelve-tone musical scale, upon which a large portion of all music (western popular music in particular) is based, may have its roots in the sound of the human voice during the course of evolution, according to a study published by the New Scientist. Analysis of recorded speech samples found peaks in acoustic energy that mirrored the distances between notes in the twelve-tone scale. Voice disorders There are many disorders that affect the human voice; these include speech impediments, and growths and lesions on the vocal folds. Talking improperly for long periods of time causes vocal loading, which is stress inflicted on the speech organs. When vocal injury is done, often an ENT specialist may be able to help, but the best treatment is the prevention of injuries through good vocal production. Voice therapy is generally delivered by a speech-language pathologist. Vocal cord nodules and polyps Vocal nodules are caused over time by repeated abuse of the vocal cords which results in soft, swollen spots on each vocal cord. These spots develop into harder, callous-like growths called nodules. The longer the abuse occurs the larger and stiffer the nodules will become. Most polyps are larger than nodules and may be called by other names, such as polypoid degeneration or Reinke's edema. Polyps are caused by a single occurrence and may require surgical removal. Irritation after the removal may then lead to nodules if additional irritation persists. Speech-language therapy teaches the patient how to eliminate the irritations permanently through habit changes and vocal hygiene. Hoarseness or breathiness that lasts for more than two weeks is a common symptom of an underlying voice disorder such as nodes or polyps and should be investigated medically.	Biology and health sciences	Human anatomy	null
38493	https://en.wikipedia.org/wiki/Genus	Genus	Genus (; : genera ) is a taxonomic rank above species and below family as used in the biological classification of living and fossil organisms as well as viruses. In binomial nomenclature, the genus name forms the first part of the binomial species name for each species within the genus. E.g. Panthera leo (lion) and Panthera onca (jaguar) are two species within the genus Panthera. Panthera is a genus within the family Felidae. The composition of a genus is determined by taxonomists. The standards for genus classification are not strictly codified, so different authorities often produce different classifications for genera. There are some general practices used, however, including the idea that a newly defined genus should fulfill these three criteria to be descriptively useful: monophyly – all descendants of an ancestral taxon are grouped together (i.e. phylogenetic analysis should clearly demonstrate both monophyly and validity as a separate lineage). reasonable compactness – a genus should not be expanded needlessly. distinctness – with respect to evolutionarily relevant criteria, i.e. ecology, morphology, or biogeography; DNA sequences are a consequence rather than a condition of diverging evolutionary lineages except in cases where they directly inhibit gene flow (e.g. postzygotic barriers). Moreover, genera should be composed of phylogenetic units of the same kind as other (analogous) genera. Etymology The term "genus" comes from Latin , a noun form cognate with ('to bear; to give birth to'). The Swedish taxonomist Carl Linnaeus popularized its use in his 1753 Species Plantarum, but the French botanist Joseph Pitton de Tournefort (1656–1708) is considered "the founder of the modern concept of genera". Use The scientific name (or the scientific epithet) of a genus is also called the generic name; in modern style guides and science, it is always capitalised. It plays a fundamental role in binomial nomenclature, the system of naming organisms, where it is combined with the scientific name of a species: see Botanical name and Specific name (zoology). Use in nomenclature The rules for the scientific names of organisms are laid down in the nomenclature codes, which allow each species a single unique name that, for animals (including protists), plants (also including algae and fungi) and prokaryotes (bacteria and archaea), is Latin and binomial in form; this contrasts with common or vernacular names, which are non-standardized, can be non-unique, and typically also vary by country and language of usage. Except for viruses, the standard format for a species name comprises the generic name, indicating the genus to which the species belongs, followed by the specific epithet, which (within that genus) is unique to the species. For example, the gray wolf's scientific name is with Canis (Latin for 'dog') being the generic name shared by the wolf's close relatives and (Latin for 'wolf') being the specific name particular to the wolf. A botanical example would be Hibiscus arnottianus, a particular species of the genus Hibiscus native to Hawaii. The specific name is written in lower-case and may be followed by subspecies names in zoology or a variety of infraspecific names in botany. When the generic name is already known from context, it may be shortened to its initial letter, for example, C. lupus in place of Canis lupus. Where species are further subdivided, the generic name (or its abbreviated form) still forms the leading portion of the scientific name, for example, for the Eurasian wolf subspecies, or as a botanical example, . Also, as visible in the above examples, the Latinised portions of the scientific names of genera and their included species (and infraspecies, where applicable) are, by convention, written in italics. The scientific names of virus species are descriptive, not binomial in form, and may or may not incorporate an indication of their containing genus; for example, the virus species "Salmonid herpesvirus 1", "Salmonid herpesvirus 2" and "Salmonid herpesvirus 3" are all within the genus Salmonivirus; however, the genus to which the species with the formal names "Everglades virus" and "Ross River virus" are assigned is Alphavirus. As with scientific names at other ranks, in all groups other than viruses, names of genera may be cited with their authorities, typically in the form "author, year" in zoology, and "standard abbreviated author name" in botany. Thus in the examples above, the genus Canis would be cited in full as "Canis Linnaeus, 1758" (zoological usage), while Hibiscus, also first established by Linnaeus but in 1753, is simply "Hibiscus L." (botanical usage). The type concept Each genus should have a designated type, although in practice there is a backlog of older names without one. In zoology, this is the type species, and the generic name is permanently associated with the type specimen of its type species. Should the specimen turn out to be assignable to another genus, the generic name linked to it becomes a junior synonym and the remaining taxa in the former genus need to be reassessed. Categories of generic name In zoological usage, taxonomic names, including those of genera, are classified as "available" or "unavailable". Available names are those published in accordance with the International Code of Zoological Nomenclature; the earliest such name for any taxon (for example, a genus) should then be selected as the "valid" (i.e., current or accepted) name for the taxon in question. Consequently, there will be more available names than valid names at any point in time; which names are currently in use depending on the judgement of taxonomists in either combining taxa described under multiple names, or splitting taxa which may bring available names previously treated as synonyms back into use. "Unavailable" names in zoology comprise names that either were not published according to the provisions of the ICZN Code, e.g., incorrect original or subsequent spellings, names published only in a thesis, and generic names published after 1930 with no type species indicated. According to "Glossary" section of the zoological Code, suppressed names (per published "Opinions" of the International Commission of Zoological Nomenclature) remain available but cannot be used as the valid name for a taxon; however, the names published in suppressed works are made unavailable via the relevant Opinion dealing with the work in question. In botany, similar concepts exist but with different labels. The botanical equivalent of zoology's "available name" is a validly published name. An invalidly published name is a or ; a rejected name is a or ; a later homonym of a validly published name is a or ; for a full list refer to the International Code of Nomenclature for algae, fungi, and plants and the work cited above by Hawksworth, 2010. In place of the "valid taxon" in zoology, the nearest equivalent in botany is "correct name" or "current name" which can, again, differ or change with alternative taxonomic treatments or new information that results in previously accepted genera being combined or split. Prokaryote and virus codes of nomenclature also exist which serve as a reference for designating currently accepted genus names as opposed to others which may be either reduced to synonymy, or, in the case of prokaryotes, relegated to a status of "names without standing in prokaryotic nomenclature". An available (zoological) or validly published (botanical) name that has been historically applied to a genus but is not regarded as the accepted (current/valid) name for the taxon is termed a synonym; some authors also include unavailable names in lists of synonyms as well as available names, such as misspellings, names previously published without fulfilling all of the requirements of the relevant nomenclatural code, and rejected or suppressed names. A particular genus name may have zero to many synonyms, the latter case generally if the genus has been known for a long time and redescribed as new by a range of subsequent workers, or if a range of genera previously considered separate taxa have subsequently been consolidated into one. For example, the World Register of Marine Species presently lists 8 genus-level synonyms for the sperm whale genus Physeter Linnaeus, 1758, and 13 for the bivalve genus Pecten O.F. Müller, 1776. Identical names (homonyms) Within the same kingdom, one generic name can apply to one genus only. However, many names have been assigned (usually unintentionally) to two or more different genera. For example, the platypus belongs to the genus Ornithorhynchus although George Shaw named it Platypus in 1799 (these two names are thus synonyms). However, the name Platypus had already been given to a group of ambrosia beetles by Johann Friedrich Wilhelm Herbst in 1793. A name that means two different things is a homonym. Since beetles and platypuses are both members of the kingdom Animalia, the name could not be used for both. Johann Friedrich Blumenbach published the replacement name Ornithorhynchus in 1800. However, a genus in one kingdom is allowed to bear a scientific name that is in use as a generic name (or the name of a taxon in another rank) in a kingdom that is governed by a different nomenclature code. Names with the same form but applying to different taxa are called "homonyms". Although this is discouraged by both the International Code of Zoological Nomenclature and the International Code of Nomenclature for algae, fungi, and plants, there are some five thousand such names in use in more than one kingdom. For instance, Anura is the name of the order of frogs but also is the name of a non-current genus of plants; Aotus is the generic name of both golden peas and night monkeys; Oenanthe is the generic name of both wheatears and water dropworts; Prunella is the generic name of both accentors and self-heal; and Proboscidea is the order of elephants and the genus of devil's claws. The name of the genus Paramecia (an extinct red alga) is also the plural of the name of the genus Paramecium (which is in the SAR supergroup), which can also lead to confusion. A list of generic homonyms (with their authorities), including both available (validly published) and selected unavailable names, has been compiled by the Interim Register of Marine and Nonmarine Genera (IRMNG). Use in higher classifications The type genus forms the base for higher taxonomic ranks, such as the family name ("Canids") based on Canis. However, this does not typically ascend more than one or two levels: the order to which dogs and wolves belong is ("Carnivores"). Numbers of accepted genera The numbers of either accepted, or all published genus names is not known precisely; Rees et al., 2020 estimate that approximately 310,000 accepted names (valid taxa) may exist, out of a total of c. 520,000 published names (including synonyms) as at end 2019, increasing at some 2,500 published generic names per year. "Official" registers of taxon names at all ranks, including genera, exist for a few groups only such as viruses and prokaryotes, while for others there are compendia with no "official" standing such as Index Fungorum for fungi, Index Nominum Algarum and AlgaeBase for algae, Index Nominum Genericorum and the International Plant Names Index for plants in general, and ferns through angiosperms, respectively, and Nomenclator Zoologicus and the Index to Organism Names for zoological names. Totals for both "all names" and estimates for "accepted names" as held in the Interim Register of Marine and Nonmarine Genera (IRMNG) are broken down further in the publication by Rees et al., 2020 cited above. The accepted names estimates are as follows, broken down by kingdom: Animalia: 239,093 accepted genus names (± 55,350) Plantae: 28,724 (± 7,721) Fungi: 10,468 (± 182) Chromista: 11,114 (± 1,268) Protozoa: 3,109 (± 1,206) Bacteria: 3,433 (± 115) Archaea: 140 (± 0) Viruses: 851 (± 0) The cited ranges of uncertainty arise because IRMNG lists "uncertain" names (not researched therein) in addition to known "accepted" names; the values quoted are the mean of "accepted" names alone (all "uncertain" names treated as unaccepted) and "accepted + uncertain" names (all "uncertain" names treated as accepted), with the associated range of uncertainty indicating these two extremes. Within Animalia, the largest phylum is Arthropoda, with 151,697 ± 33,160 accepted genus names, of which 114,387 ± 27,654 are insects (class Insecta). Within Plantae, Tracheophyta (vascular plants) make up the largest component, with 23,236 ± 5,379 accepted genus names, of which 20,845 ± 4,494 are angiosperms (superclass Angiospermae). By comparison, the 2018 annual edition of the Catalogue of Life (estimated >90% complete, for extant species in the main) contains currently 175,363 "accepted" genus names for 1,744,204 living and 59,284 extinct species, also including genus names only (no species) for some groups. Genus size The number of species in genera varies considerably among taxonomic groups. For instance, among (non-avian) reptiles, which have about 1180 genera, the most (>300) have only 1 species, ~360 have between 2 and 4 species, 260 have 5–10 species, ~200 have 11–50 species, and only 27 genera have more than 50 species. However, some insect genera such as the bee genera Lasioglossum and Andrena have over 1000 species each. The largest flowering plant genus, Astragalus, contains over 3,000 species. Which species are assigned to a genus is somewhat arbitrary. Although all species within a genus are supposed to be "similar", there are no objective criteria for grouping species into genera. There is much debate among zoologists about whether enormous, species-rich genera should be maintained, as it is extremely difficult to come up with identification keys or even character sets that distinguish all species. Hence, many taxonomists argue in favor of breaking down large genera. For instance, the lizard genus Anolis has been suggested to be broken down into 8 or so different genera which would bring its ~400 species to smaller, more manageable subsets.	Biology and health sciences	Taxonomic rank	Biology
38498	https://en.wikipedia.org/wiki/Frog	Frog	A frog is any member of a diverse and largely carnivorous group of short-bodied, tailless amphibians composing the order Anura (coming from the Ancient Greek , literally 'without tail'). The oldest fossil "proto-frog" Triadobatrachus is known from the Early Triassic of Madagascar (250million years ago), but molecular clock dating suggests their split from other amphibians may extend further back to the Permian, 265million years ago. Frogs are widely distributed, ranging from the tropics to subarctic regions, but the greatest concentration of species diversity is in tropical rainforest. Frogs account for around 88% of extant amphibian species. They are also one of the five most diverse vertebrate orders. Warty frog species tend to be called toads, but the distinction between frogs and toads is informal, not from taxonomy or evolutionary history. An adult frog has a stout body, protruding eyes, anteriorly-attached tongue, limbs folded underneath, and no tail (the tail of tailed frogs is an extension of the male cloaca). Frogs have glandular skin, with secretions ranging from distasteful to toxic. Their skin varies in colour from well-camouflaged dappled brown, grey and green to vivid patterns of bright red or yellow and black to show toxicity and ward off predators. Adult frogs live in fresh water and on dry land; some species are adapted for living underground or in trees. Frogs typically lay their eggs in the water. The eggs hatch into aquatic larvae called tadpoles that have tails and internal gills. They have highly specialised rasping mouth parts suitable for herbivorous, omnivorous or planktivorous diets. The life cycle is completed when they metamorphose into adults. A few species deposit eggs on land or bypass the tadpole stage. Adult frogs generally have a carnivorous diet consisting of small invertebrates, but omnivorous species exist and a few feed on plant matter. Frog skin has a rich microbiome which is important to their health. Frogs are extremely efficient at converting what they eat into body mass. They are an important food source for predators and part of the food web dynamics of many of the world's ecosystems. The skin is semi-permeable, making them susceptible to dehydration, so they either live in moist places or have special adaptations to deal with dry habitats. Frogs produce a wide range of vocalisations, particularly in their breeding season, and exhibit many different kinds of complex behaviors to attract mates, to fend off predators and to generally survive. Frogs are valued as food by humans and also have many cultural roles in literature, symbolism and religion. They are also seen as environmental bellwethers, with declines in frog populations often viewed as early warning signs of environmental damage. Frog populations have declined significantly since the 1950s. More than one third of species are considered to be threatened with extinction and over 120 are believed to have become extinct since the 1980s. The number of malformations among frogs is on the rise and an emerging fungal disease, chytridiomycosis, has spread around the world. Conservation biologists are working to understand the causes of these problems and to resolve them. Etymology and taxonomy The use of the common names frog and toad has no taxonomic justification. From a classification perspective, all members of the order Anura are frogs, but only members of the family Bufonidae are considered "true toads". The use of the term frog in common names usually refers to species that are aquatic or semi-aquatic and have smooth, moist skins; the term toad generally refers to species that are terrestrial with dry, warty skins. There are numerous exceptions to this rule. The European fire-bellied toad (Bombina bombina) has a slightly warty skin and prefers a watery habitat whereas the Panamanian golden frog (Atelopus zeteki) is in the toad family Bufonidae and has a smooth skin. Etymology The origin of the order name Anura—and its original spelling Anoures—is the Ancient Greek alpha privative prefix ( from before a vowel) 'without', and () 'animal tail'. meaning "tailless". It refers to the tailless character of these amphibians. The origins of the word frog are uncertain and debated. The word is first attested in Old English as , but the usual Old English word for the frog was (with variants such as and ), and it is agreed that the word frog is somehow related to this. Old English remained in dialectal use in English as frosh and frosk into the nineteenth century, and is paralleled widely in other Germanic languages, with examples in the modern languages including German , Norwegian , Icelandic , and Dutch . These words allow reconstruction of a Common Germanic ancestor . The third edition of the Oxford English Dictionary finds that the etymology of is uncertain, but agrees with arguments that it could plausibly derive from a Proto-Indo-European base along the lines of , meaning 'jump'. How Old English gave rise to is, however, uncertain, as the development does not involve a regular sound-change. Instead, it seems that there was a trend in Old English to coin nicknames for animals ending in -g, with examples—themselves all of uncertain etymology—including dog, hog, pig, stag, and . Frog appears to have been adapted from as part of this trend. Meanwhile, the word toad, first attested as Old English , is unique to English and is likewise of uncertain etymology. It is the basis for the word tadpole, first attested as Middle English , apparently meaning 'toad-head'. Taxonomy About 88% of amphibian species are classified in the order Anura. These include over 7,700 species in 59 families, of which the Hylidae (1062 spp.), Strabomantidae (807 spp.), Microhylidae (758 spp.), and Bufonidae (657 spp.) are the richest in species. The Anura include all modern frogs and any fossil species that fit within the anuran definition. The characteristics of anuran adults include: 9 or fewer presacral vertebrae, the presence of a urostyle formed of fused vertebrae, no tail, a long and forward-sloping ilium, shorter fore limbs than hind limbs, radius and ulna fused, tibia and fibula fused, elongated ankle bones, absence of a prefrontal bone, presence of a hyoid plate, a lower jaw without teeth (with the exception of Gastrotheca guentheri) consisting of three pairs of bones (angulosplenial, dentary, and mentomeckelian, with the last pair being absent in Pipoidea), an unsupported tongue, lymph spaces underneath the skin, and a muscle, the protractor lentis, attached to the lens of the eye. The anuran larva or tadpole has a single central respiratory spiracle and mouthparts consisting of keratinous beaks and denticles. Frogs and toads are broadly classified into three suborders: Archaeobatrachia, which includes four families of primitive frogs; Mesobatrachia, which includes five families of more evolutionary intermediate frogs; and Neobatrachia, by far the largest group, which contains the remaining families of modern frogs, including most common species throughout the world. The suborder Neobatrachia is further divided into the two superfamilies Hyloidea and Ranoidea. This classification is based on such morphological features as the number of vertebrae, the structure of the pectoral girdle, and the morphology of tadpoles. While this classification is largely accepted, relationships among families of frogs are still debated. Some species of anurans hybridise readily. For instance, the edible frog (Pelophylax esculentus) is a hybrid between the pool frog (P. lessonae) and the marsh frog (P. ridibundus). The fire-bellied toads Bombina bombina and B. variegata are similar in forming hybrids. These are less fertile than their parents, giving rise to a hybrid zone where the hybrids are prevalent. Evolution The origins and evolutionary relationships between the three main groups of amphibians are hotly debated. A molecular phylogeny based on rDNA analysis dating from 2005 suggests that salamanders and caecilians are more closely related to each other than they are to frogs and the divergence of the three groups took place in the Paleozoic or early Mesozoic before the break-up of the supercontinent Pangaea and soon after their divergence from the lobe-finned fishes. This would help account for the relative scarcity of amphibian fossils from the period before the groups split. Another molecular phylogenetic analysis conducted about the same time concluded that lissamphibians first appeared about 330 million years ago and that the temnospondyl-origin hypothesis is more credible than other theories. The neobatrachians seemed to have originated in Africa/India, the salamanders in East Asia and the caecilians in tropical Pangaea. Other researchers, while agreeing with the main thrust of this study, questioned the choice of calibration points used to synchronise the data. They proposed that the date of lissamphibian diversification should be placed in the Permian, rather less than 300 million years ago, a date in better agreement with the palaeontological data. A further study in 2011 using both extinct and living taxa sampled for morphological, as well as molecular data, came to the conclusion that Lissamphibia is monophyletic and that it should be nested within Lepospondyli rather than within Temnospondyli. The study postulated that Lissamphibia originated no earlier than the late Carboniferous, some 290 to 305 million years ago. The split between Anura and Caudata was estimated as taking place 292 million years ago, rather later than most molecular studies suggest, with the caecilians splitting off 239 million years ago. In 2008, Gerobatrachus hottoni, a temnospondyl with many frog- and salamander-like characteristics, was discovered in Texas. It dated back 290 million years and was hailed as a missing link, a stem batrachian close to the common ancestor of frogs and salamanders, consistent with the widely accepted hypothesis that frogs and salamanders are more closely related to each other (forming a clade called Batrachia) than they are to caecilians. However, others have suggested that Gerobatrachus hottoni was only a dissorophoid temnospondyl unrelated to extant amphibians. Salientia (Latin salire (salio), "to jump") is the name of the total group that includes modern frogs in the order Anura as well as their close fossil relatives, the "proto-frogs" or "stem-frogs". The common features possessed by these proto-frogs include 14 presacral vertebrae (modern frogs have eight or 9), a long and forward-sloping ilium in the pelvis, the presence of a frontoparietal bone, and a lower jaw without teeth. The earliest known amphibians that were more closely related to frogs than to salamanders are Triadobatrachus massinoti, from the early Triassic period of Madagascar (about 250 million years ago), and Czatkobatrachus polonicus, from the Early Triassic of Poland (about the same age as Triadobatrachus). The skull of Triadobatrachus is frog-like, being broad with large eye sockets, but the fossil has features diverging from modern frogs. These include a longer body with more vertebrae. The tail has separate vertebrae unlike the fused urostyle or coccyx in modern frogs. The tibia and fibula bones are also separate, making it probable that Triadobatrachus was not an efficient leaper. A 2019 study has noted the presence of Salientia from the Chinle Formation, and suggested that anurans might have first appeared during the Late Triassic. On the basis of fossil evidence, the earliest known "true frogs" that fall into the anuran lineage proper all lived in the early Jurassic period. One such early frog species, Prosalirus bitis, was discovered in 1995 in the Kayenta Formation of Arizona and dates back to the Early Jurassic epoch (199.6 to 175 million years ago), making Prosalirus somewhat more recent than Triadobatrachus. Like the latter, Prosalirus did not have greatly enlarged legs, but had the typical three-pronged pelvic structure of modern frogs. Unlike Triadobatrachus, Prosalirus had already lost nearly all of its tail and was well adapted for jumping. Another Early Jurassic frog is Vieraella herbsti, which is known only from dorsal and ventral impressions of a single animal and was estimated to be from snout to vent. Notobatrachus degiustoi from the middle Jurassic is slightly younger, about 155–170 million years old. The main evolutionary changes in this species involved the shortening of the body and the loss of the tail. Tadpoles of N. degiustoi constitute the oldest tadpoles found as of 2024, dating back to 168–161 million years ago. These tadpoles also showed adaptations for filter-feeding, implying residence in temporary pools by filter-feeding larvae was already commonplace. The evolution of modern Anura likely was complete by the Jurassic period. Since then, evolutionary changes in chromosome numbers have taken place about 20 times faster in mammals than in frogs, which means speciation is occurring more rapidly in mammals. According to genetic studies, the families Hyloidea, Microhylidae, and the clade Natatanura (comprising about 88% of living frogs) diversified simultaneously some 66 million years ago, soon after the Cretaceous–Paleogene extinction event associated with the Chicxulub impactor. All origins of arboreality (e.g. in Hyloidea and Natatanura) follow from that time and the resurgence of forest that occurred afterwards. Frog fossils have been found on all of the Earth's continents. In 2020, it was announced that 40 million year old helmeted frog fossils had been discovered by a team of vertebrate palaeontologists in Seymour Island on the Antarctic Peninsula, indicating that this region was once home to frogs related to those now living in South American Nothofagus forest. Phylogeny A cladogram showing the relationships of the different families of frogs in the clade Anura can be seen in the table below. This diagram, in the form of a tree, shows how each frog family is related to other families, with each node representing a point of common ancestry. It is based on Frost et al. (2006), Heinicke et al. (2009) and Pyron and Wiens (2011). Morphology and physiology Frogs have no tail, except as larvae, and most have long hind legs, elongated ankle bones, webbed toes, no claws, large eyes, and a smooth or warty skin. They have short vertebral columns, with no more than 10 free vertebrae and fused tailbones (urostyle or coccyx). Frogs range in size from Paedophryne amauensis of Papua New Guinea that is in snout–vent length to the up to about and goliath frog (Conraua goliath) of central Africa. There are prehistoric, extinct species that reached even larger sizes. Feet and legs The structure of the feet and legs varies greatly among frog species, depending in part on whether they live primarily on the ground, in water, in trees, or in burrows. Adult anurans have four fingers on the hands and five toes on the feet, but the smallest species often have hands and feet where some of the digits are vestigial. Frogs must be able to move quickly through their environment to catch prey and escape predators, and numerous adaptations help them to do so. Most frogs are either proficient at jumping or are descended from ancestors that were, with much of the musculoskeletal morphology modified for this purpose. The tibia, fibula, and tarsals have been fused into a single strong bone, as have the radius and ulna in the fore limbs (which must absorb the impact on landing). The metatarsals have become elongated to add to the leg length and allow frogs to push against the ground for a longer period on take-off. The ilium has elongated and formed a mobile joint with the sacrum which, in specialist jumpers such as ranids and hylids, functions as an additional limb joint to further power the leaps. The tail vertebrae have fused into a urostyle which is retracted inside the pelvis. This enables the force to be transferred from the legs to the body during a leap. The muscular system has been similarly modified. The hind limbs of ancestral frogs presumably contained pairs of muscles which would act in opposition (one muscle to flex the knee, a different muscle to extend it), as is seen in most other limbed animals. However, in modern frogs, almost all muscles have been modified to contribute to the action of jumping, with only a few small muscles remaining to bring the limb back to the starting position and maintain posture. The muscles have also been greatly enlarged, with the main leg muscles accounting for over 17% of the total mass of frogs. Many frogs have webbed feet and the degree of webbing is directly proportional to the amount of time the species spends in the water. The completely aquatic African dwarf frog (Hymenochirus sp.) has fully webbed toes, whereas those of White's tree frog (Litoria caerulea), an arboreal species, are only a quarter or half webbed. Exceptions include flying frogs in the Hylidae and Rhacophoridae, which also have fully webbed toes used in gliding. Arboreal frogs have pads located on the ends of their toes to help grip vertical surfaces. These are not suction pads, the surface consisting instead of columnar cells with flat tops with small gaps between them lubricated by mucous glands. When the frog applies pressure, the cells adhere to irregularities on the surface and the grip is maintained through surface tension. This allows the frog to climb on smooth surfaces, but the system does not function efficiently when the pads are excessively wet. In many arboreal frogs, a small "intercalary structure" on each toe increases the surface area touching the substrate. Furthermore, many arboreal frogs have hip joints that allow both hopping and walking. Some frogs that live high in trees even possess an elaborate degree of webbing between their toes. This allows the frogs to "parachute" or make a controlled glide from one position in the canopy to another. Ground-dwelling frogs generally lack the adaptations of aquatic and arboreal frogs. Most have smaller toe pads, if any, and little webbing. Some burrowing frogs such as Couch's spadefoot (Scaphiopus couchii) have a flap-like toe extension on the hind feet, a keratinised tubercle often referred to as a spade, that helps them to burrow. Sometimes during the tadpole stage, one of the developing rear legs is eaten by a predator such as a dragonfly nymph. In some cases, the full leg still grows, but in others it does not, although the frog may still live out its normal lifespan with only three limbs. Occasionally, a parasitic flatworm (Ribeiroia ondatrae) digs into the rear of a tadpole, causing a rearrangement of the limb bud cells and the frog develops one or more extra legs. Skin A frog's skin is protective, has a respiratory function, can absorb water, and helps control body temperature. It has many glands, particularly on the head and back, which often exude distasteful and toxic substances (granular glands). The secretion is often sticky and helps keep the skin moist, protects against the entry of moulds and bacteria, and makes the animal slippery and more able to escape from predators. The skin is shed every few weeks. It usually splits down the middle of the back and across the belly, and the frog pulls its arms and legs free. The sloughed skin is then worked towards the head where it is quickly eaten. Being cold-blooded, frogs have to adopt suitable behaviour patterns to regulate their temperature. To warm up, they can move into the sun or onto a warm surface; if they overheat, they can move into the shade or adopt a stance that exposes the minimum area of skin to the air. This posture is also used to prevent water loss and involves the frog squatting close to the substrate with its hands and feet tucked under its chin and body. The colour of a frog's skin is used for thermoregulation. In cool damp conditions, the colour will be darker than on a hot dry day. The grey foam-nest tree frog (Chiromantis xerampelina) is even able to turn white to minimise the chance of overheating. Many frogs are able to absorb water and oxygen directly through the skin, especially around the pelvic area, but the permeability of a frog's skin can also result in water loss. Glands located all over the body exude mucus which helps keep the skin moist and reduces evaporation. Some glands on the hands and chest of males are specialised to produce sticky secretions to aid in amplexus. Similar glands in tree frogs produce a glue-like substance on the adhesive discs of the feet. Some arboreal frogs reduce water loss by having a waterproof layer of skin, and several South American species coat their skin with a waxy secretion. Other frogs have adopted behaviours to conserve water, including becoming nocturnal and resting in a water-conserving position. Some frogs may also rest in large groups with each frog pressed against its neighbours. This reduces the amount of skin exposed to the air or a dry surface, and thus reduces water loss. Woodhouse's toad (Bufo woodhousii), if given access to water after confinement in a dry location, sits in the shallows to rehydrate. The male hairy frog (Trichobatrachus robustus) has dermal papillae projecting from its lower back and thighs, giving it a bristly appearance. These contain blood vessels and are thought to increase the area of the skin available for respiration. Some species have bony plates embedded in the skin, a trait that appears to have evolved independently several times. In certain other species, the skin at the top of the head is compacted and the connective tissue of the dermis is co-ossified with the bones of the skull (exostosis). Camouflage is a common defensive mechanism in frogs. Features such as warts and skin folds are usually on ground-dwelling frogs, for whom smooth skin would not provide such effective camouflage. Certain frogs change colour between night and day, as light and moisture stimulate the pigment cells and cause them to expand or contract. Some are even able to control their skin texture. The Pacific tree frog (Pseudacris regilla) has green and brown morphs, plain or spotted, and changes colour depending on the time of year and general background colour. The Wood frog (Lithobates sylvaticus) uses disruptive coloration including black eye markings similar to voids between leaves, bands of the dorsal skin (dorsolateral dermal plica) similar to a leaf midrib as well as stains, spots and leg stripes similar to fallen leaf features. Respiration and circulation Like other amphibians, oxygen can pass through their highly permeable skins. This unique feature allows them to remain in places without access to the air, respiring through their skins. Ribs are generally absent, so the lungs are filled by buccal pumping and a frog deprived of its lungs can maintain its body functions without them. The fully aquatic Bornean flat-headed frog (Barbourula kalimantanensis) is the first frog known to lack lungs entirely. Frogs have three-chambered hearts, a feature they share with lizards. Oxygenated blood from the lungs and de-oxygenated blood from the respiring tissues enter the heart through separate atria. When these chambers contract, the two blood streams pass into a common ventricle before being pumped via a spiral valve to the appropriate vessel, the aorta for oxygenated blood and pulmonary artery for deoxygenated blood. Some species of frog have adaptations that allow them to survive in oxygen deficient water. The Titicaca water frog (Telmatobius culeus) is one such species and has wrinkly skin that increases its surface area to enhance gas exchange. It normally makes no use of its rudimentary lungs but will sometimes raise and lower its body rhythmically while on the lake bed to increase the flow of water around it. Digestion and excretion Frogs have maxillary teeth along their upper jaw which are used to hold food before it is swallowed. These teeth are very weak, and cannot be used to chew or catch and harm agile prey. Instead, the frog uses its sticky, cleft tongue to catch insects and other small moving prey. The tongue normally lies coiled in the mouth, free at the back and attached to the mandible at the front. It can be shot out and retracted at great speed. In amphibians there are salvary glands on the tongue, which in frogs produce what is called a two-phase viscoelastic fluid. When exposed to pressure, like when the tongue is wrapping around a prey, it becomes runny and covers the prey's body. As the pressure drops, it returns to a thick and elastic state, which gives the tongue an extra grip. Some frogs have no tongue and just stuff food into their mouths with their hands. The African bullfrog (Pyxicephalus), which preys on relatively large animals such as mice and other frogs, has cone shaped bony projections called odontoid processes at the front of the lower jaw which function like teeth. The eyes assist in the swallowing of food as they can be retracted through holes in the skull and help push food down the throat. The food then moves through the oesophagus into the stomach where digestive enzymes are added and it is churned up. It then proceeds to the small intestine (duodenum and ileum) where most digestion occurs. Pancreatic juice from the pancreas, and bile, produced by the liver and stored in the gallbladder, are secreted into the small intestine, where the fluids digest the food and the nutrients are absorbed. The food residue passes into the large intestine where excess water is removed and the wastes are passed out through the cloaca. Although adapted to terrestrial life, frogs resemble freshwater fish in their inability to conserve body water effectively. When they are on land, much water is lost by evaporation from the skin. The excretory system is similar to that of mammals and there are two kidneys that remove nitrogenous products from the blood. Frogs produce large quantities of dilute urine in order to flush out toxic products from the kidney tubules. The nitrogen is excreted as ammonia by tadpoles and aquatic frogs but mainly as urea, a less toxic product, by most terrestrial adults. A few species of tree frog with little access to water excrete the even less toxic uric acid. The urine passes along paired ureters to the urinary bladder from which it is vented periodically into the cloaca. All bodily wastes exit the body through the cloaca which terminates in a cloacal vent. Reproductive system In the male frog, the two testes are attached to the kidneys and semen passes into the kidneys through fine tubes called efferent ducts. It then travels on through the ureters, which are consequently known as urinogenital ducts. There is no penis, and sperm is ejected from the cloaca directly onto the eggs as the female lays them. The ovaries of the female frog are beside the kidneys and the eggs pass down a pair of oviducts and through the cloaca to the exterior. When frogs mate, the male climbs on the back of the female and wraps his fore limbs round her body, either behind the front legs or just in front of the hind legs. This position is called amplexus and may be held for several days. The male frog has certain hormone-dependent secondary sexual characteristics. These include the development of special pads on his thumbs in the breeding season, to give him a firm hold. The grip of the male frog during amplexus stimulates the female to release eggs, usually wrapped in jelly, as spawn. In many species the male is smaller and slimmer than the female. Males have vocal cords and make a range of croaks, particularly in the breeding season, and in some species they also have vocal sacs to amplify the sound. Nervous system Frogs have a highly developed nervous system that consists of a brain, spinal cord and nerves. Many parts of frog brains correspond with those of humans. It consists of two olfactory lobes, two cerebral hemispheres, a pineal body, two optic lobes, a cerebellum and a medulla oblongata. Muscular coordination and posture are controlled by the cerebellum, and the medulla oblongata regulates respiration, digestion and other automatic functions. The relative size of the cerebrum in frogs is much smaller than it is in humans. Frogs have ten pairs of cranial nerves which pass information from the outside directly to the brain, and ten pairs of spinal nerves which pass information from the extremities to the brain through the spinal cord. By contrast, all amniotes (mammals, birds and reptiles) have twelve pairs of cranial nerves. Sight The eyes of most frogs are located on either side of the head near the top and project outwards as hemispherical bulges. They provide binocular vision over a field of 100° to the front and a total visual field of almost 360°. They may be the only part of an otherwise submerged frog to protrude from the water. Each eye has closable upper and lower lids and a nictitating membrane which provides further protection, especially when the frog is swimming. Members of the aquatic family Pipidae have the eyes located at the top of the head, a position better suited for detecting prey in the water above. The irises come in a range of colours and the pupils in a range of shapes. The common toad (Bufo bufo) has golden irises and horizontal slit-like pupils, the red-eyed tree frog (Agalychnis callidryas) has vertical slit pupils, the poison dart frog has dark irises, the fire-bellied toad (Bombina spp.) has triangular pupils and the tomato frog (Dyscophus spp.) has circular ones. The irises of the southern toad (Anaxyrus terrestris) are patterned so as to blend in with the surrounding camouflaged skin. The distant vision of a frog is better than its near vision. Calling frogs will quickly become silent when they see an intruder or even a moving shadow but the closer an object is, the less well it is seen. When a frog shoots out its tongue to catch an insect it is reacting to a small moving object that it cannot see well and must line it up precisely beforehand because it shuts its eyes as the tongue is extended. Although it was formerly debated, more recent research has shown that frogs can see in colour, even in very low light. Hearing Frogs can hear both in the air and below water. They do not have external ears; the eardrums (tympanic membranes) are directly exposed or may be covered by a layer of skin and are visible as a circular area just behind the eye. The size and distance apart of the eardrums is related to the frequency and wavelength at which the frog calls. In some species such as the bullfrog, the size of the tympanum indicates the sex of the frog; males have tympani that are larger than their eyes while in females, the eyes and tympani are much the same size. A noise causes the tympanum to vibrate and the sound is transmitted to the middle and inner ear. The middle ear contains semicircular canals which help control balance and orientation. In the inner ear, the auditory hair cells are arranged in two areas of the cochlea, the basilar papilla and the amphibian papilla. The former detects high frequencies and the latter low frequencies. Because the cochlea is short, frogs use electrical tuning to extend their range of audible frequencies and help discriminate different sounds. This arrangement enables detection of the territorial and breeding calls of their conspecifics. In some species that inhabit arid regions, the sound of thunder or heavy rain may arouse them from a dormant state. A frog may be startled by an unexpected noise but it will not usually take any action until it has located the source of the sound by sight. Call The call or croak of a frog is unique to its species. Frogs create this sound by passing air through the larynx in the throat. In most calling frogs, the sound is amplified by one or more vocal sacs, membranes of skin under the throat or on the corner of the mouth, that distend during the amplification of the call. Some frog calls are so loud that they can be heard up to a mile (1.6km) away. Additionally, some species have been found to use man-made structures such as drain pipes for artificial amplification of their call. The coastal tailed frog (Ascaphus truei) lives in mountain streams in North America and does not vocalise. The main function of calling is for male frogs to attract mates. Males may call individually or there may be a chorus of sound where numerous males have converged on breeding sites. In many frog species, such as the common tree frog (Polypedates leucomystax), females reply to males' calls, which acts to reinforce reproductive activity in a breeding colony. Female frogs prefer males that produce sounds of greater intensity and lower frequency, attributes that stand out in a crowd. The rationale for this is thought to be that by demonstrating his prowess, the male shows his fitness to produce superior offspring. A different call is emitted by a male frog or unreceptive female when mounted by another male. This is a distinct chirruping sound and is accompanied by a vibration of the body. Tree frogs and some non-aquatic species have a rain call that they make on the basis of humidity cues prior to a shower. Many species also have a territorial call that is used to drive away other males. All of these calls are emitted with the mouth of the frog closed. A distress call, emitted by some frogs when they are in danger, is produced with the mouth open resulting in a higher-pitched call. It is typically used when the frog has been grabbed by a predator and may serve to distract or disorient the attacker so that it releases the frog. Many species of frog have deep calls. The croak of the American bullfrog (Rana catesbiana) is sometimes written as "jug o' rum". The Pacific tree frog (Pseudacris regilla) produces the onomatopoeic "ribbit" often heard in films. Other renderings of frog calls into speech include "brekekekex koax koax", the call of the marsh frog (Pelophylax ridibundus) in The Frogs, an Ancient Greek comic drama by Aristophanes. The calls of the Concave-eared torrent frog (Amolops tormotus) are unusual in many aspects. The males are notable for their varieties of calls where upward and downward frequency modulations take place. When they communicate, they produce calls that fall in the ultrasound frequency range. The last aspect that makes this species of frog's calls unusual is that nonlinear acoustic phenomena are important components in their acoustic signals. Torpor During extreme conditions, some frogs enter a state of torpor and remain inactive for months. In colder regions, many species of frog hibernate in winter. Those that live on land such as the American toad (Bufo americanus) dig a burrow and make a hibernaculum in which to lie dormant. Others, less proficient at digging, find a crevice or bury themselves in dead leaves. Aquatic species such as the American bullfrog (Rana catesbeiana) normally sink to the bottom of the pond where they lie, semi-immersed in mud but still able to access the oxygen dissolved in the water. Their metabolism slows down and they live on their energy reserves. Some frogs such as the wood frog, moor frog, or spring peeper can even survive being frozen. Ice crystals form under the skin and in the body cavity but the essential organs are protected from freezing by a high concentration of glucose. An apparently lifeless, frozen frog can resume respiration and its heartbeat can restart when conditions warm up. At the other extreme, the striped burrowing frog (Cyclorana alboguttata) regularly aestivates during the hot, dry season in Australia, surviving in a dormant state without access to food and water for nine or ten months of the year. It burrows underground and curls up inside a protective cocoon formed by its shed skin. Researchers at the University of Queensland have found that during aestivation, the metabolism of the frog is altered and the operational efficiency of the mitochondria is increased. This means that the limited amount of energy available to the comatose frog is used in a more efficient manner. This survival mechanism is only useful to animals that remain completely unconscious for an extended period of time and whose energy requirements are low because they are cold-blooded and have no need to generate heat. Other research showed that, to provide these energy requirements, muscles atrophy, but hind limb muscles are preferentially unaffected. Frogs have been found to have upper critical temperatures of around 41 degrees Celsius. Locomotion Different species of frog use a number of methods of moving around including jumping, running, walking, swimming, burrowing, climbing and gliding. Jumping Frogs are generally recognised as exceptional jumpers and, relative to their size, the best jumpers of all vertebrates. The striped rocket frog, Litoria nasuta, can leap over , a distance that is more than fifty times its body length of . There are tremendous differences between species in jumping capability. Within a species, jump distance increases with increasing size, but relative jumping distance (body-lengths jumped) decreases. The Indian skipper frog (Euphlyctis cyanophlyctis) has the ability to leap out of the water from a position floating on the surface. The tiny northern cricket frog (Acris crepitans) can "skitter" across the surface of a pond with a series of short rapid jumps. Slow-motion photography shows that the muscles have passive flexibility. They are first stretched while the frog is still in the crouched position, then they are contracted before being stretched again to launch the frog into the air. The fore legs are folded against the chest and the hind legs remain in the extended, streamlined position for the duration of the jump. In some extremely capable jumpers, such as the Cuban tree frog (Osteopilus septentrionalis) and the northern leopard frog (Rana pipiens), the peak power exerted during a jump can exceed that which the muscle is theoretically capable of producing. When the muscles contract, the energy is first transferred into the stretched tendon which is wrapped around the ankle bone. Then the muscles stretch again at the same time as the tendon releases its energy like a catapult to produce a powerful acceleration beyond the limits of muscle-powered acceleration. A similar mechanism has been documented in locusts and grasshoppers. Early hatching of froglets can have negative effects on frog jumping performance and overall locomotion. The hindlimbs are unable to completely form, which results in them being shorter and much weaker relative to a normal hatching froglet. Early hatching froglets may tend to depend on other forms of locomotion more often, such as swimming and walking. Walking and running Frogs in the families Bufonidae, Rhinophrynidae, and Microhylidae have short back legs and tend to walk rather than jump. When they try to move rapidly, they speed up the rate of movement of their limbs or resort to an ungainly hopping gait. The Great Plains narrow-mouthed toad (Gastrophryne olivacea) has been described as having a gait that is "a combination of running and short hops that are usually only an inch or two in length". In an experiment, Fowler's toad (Bufo fowleri) was placed on a treadmill which was turned at varying speeds. By measuring the toad's uptake of oxygen it was found that hopping was an inefficient use of resources during sustained locomotion but was a useful strategy during short bursts of high-intensity activity. The red-legged running frog (Kassina maculata) has short, slim hind limbs unsuited to jumping. It can move fast by using a running gait in which the two hind legs are used alternately. Slow-motion photography shows, unlike a horse that can trot or gallop, the frog's gait remained similar at slow, medium, and fast speeds. This species can also climb trees and shrubs, and does so at night to catch insects. The Indian skipper frog (Euphlyctis cyanophlyctis) has broad feet and can run across the surface of the water for several metres (yards). Swimming Frogs that live in or visit water have adaptations that improve their swimming abilities. The hind limbs are heavily muscled and strong. The webbing between the toes of the hind feet increases the area of the foot and helps propel the frog powerfully through the water. Members of the family Pipidae are wholly aquatic and show the most marked specialisation. They have inflexible vertebral columns, flattened, streamlined bodies, lateral line systems, and powerful hind limbs with large webbed feet. Tadpoles mostly have large tail fins which provide thrust when the tail is moved from side to side. Burrowing Some frogs have become adapted for burrowing and a life underground. They tend to have rounded bodies, short limbs, small heads with bulging eyes, and hind feet adapted for excavation. An extreme example of this is the purple frog (Nasikabatrachus sahyadrensis) from southern India which feeds on termites and spends almost its whole life underground. It emerges briefly during the monsoon to mate and breed in temporary pools. It has a tiny head with a pointed snout and a plump, rounded body. Because of this fossorial existence, it was first described in 2003, being new to the scientific community at that time, although previously known to local people. The spadefoot toads of North America are also adapted to underground life. The Plains spadefoot toad (Spea bombifrons) is typical and has a flap of keratinised bone attached to one of the metatarsals of the hind feet which it uses to dig itself backwards into the ground. As it digs, the toad wriggles its hips from side to side to sink into the loose soil. It has a shallow burrow in the summer from which it emerges at night to forage. In winter, it digs much deeper and has been recorded at a depth of . The tunnel is filled with soil and the toad hibernates in a small chamber at the end. During this time, urea accumulates in its tissues and water is drawn in from the surrounding damp soil by osmosis to supply the toad's needs. Spadefoot toads are "explosive breeders", all emerging from their burrows at the same time and converging on temporary pools, attracted to one of these by the calling of the first male to find a suitable breeding location. The burrowing frogs of Australia have a rather different lifestyle. The western spotted frog (Heleioporus albopunctatus) digs a burrow beside a river or in the bed of an ephemeral stream and regularly emerges to forage. Mating takes place and eggs are laid in a foam nest inside the burrow. The eggs partially develop there, but do not hatch until they are submerged following heavy rainfall. The tadpoles then swim out into the open water and rapidly complete their development. Madagascan burrowing frogs are less fossorial and mostly bury themselves in leaf litter. One of these, the green burrowing frog (Scaphiophryne marmorata), has a flattened head with a short snout and well-developed metatarsal tubercles on its hind feet to help with excavation. It also has greatly enlarged terminal discs on its fore feet that help it to clamber around in bushes. It breeds in temporary pools that form after rains. Climbing Tree frogs live high in the canopy, where they scramble around on the branches, twigs, and leaves, sometimes never coming down to earth. The "true" tree frogs belong to the family Hylidae, but members of other frog families have independently adopted an arboreal habit, a case of convergent evolution. These include the glass frogs (Centrolenidae), the bush frogs (Hyperoliidae), some of the narrow-mouthed frogs (Microhylidae), and the shrub frogs (Rhacophoridae). Most tree frogs are under in length, with long legs and long toes with adhesive pads on the tips. The surface of the toe pads is formed from a closely packed layer of flat-topped, hexagonal epidermal cells separated by grooves into which glands secrete mucus. These toe pads, moistened by the mucus, provide the grip on any wet or dry surface, including glass. The forces involved include boundary friction of the toe pad epidermis on the surface and also surface tension and viscosity. Tree frogs are very acrobatic and can catch insects while hanging by one toe from a twig or clutching onto the blade of a windswept reed. Some members of the subfamily Phyllomedusinae have opposable toes on their feet. The reticulated leaf frog (Phyllomedusa ayeaye) has a single opposed digit on each fore foot and two opposed digits on its hind feet. This allows it to grasp the stems of bushes as it clambers around in its riverside habitat. Gliding During the evolutionary history of frogs, several different groups have independently taken to the air. Some frogs in the tropical rainforest are specially adapted for gliding from tree to tree or parachuting to the forest floor. Typical of them is Wallace's flying frog (Rhacophorus nigropalmatus) from Malaysia and Borneo. It has large feet with the fingertips expanded into flat adhesive discs and the digits fully webbed. Flaps of skin occur on the lateral margins of the limbs and across the tail region. With the digits splayed, the limbs outstretched, and these flaps spread, it can glide considerable distances, but is unable to undertake powered flight. It can alter its direction of travel and navigate distances of up to between trees. Life history Reproduction Two main types of reproduction occur in frogs, prolonged breeding and explosive breeding. In the former, adopted by the majority of species, adult frogs at certain times of year assemble at a pond, lake or stream to breed. Many frogs return to the bodies of water in which they developed as larvae. This often results in annual migrations involving thousands of individuals. In explosive breeders, mature adult frogs arrive at breeding sites in response to certain trigger factors such as rainfall occurring in an arid area. In these frogs, mating and spawning take place promptly and the speed of larval growth is rapid in order to make use of the ephemeral pools before they dry up. Among prolonged breeders, males usually arrive at the breeding site first and remain there for some time whereas females tend to arrive later and depart soon after they have spawned. This means that males outnumber females at the water's edge and defend territories from which they expel other males. They advertise their presence by calling, often alternating their croaks with neighbouring frogs. Larger, stronger males tend to have deeper calls and maintain higher quality territories. Females select their mates at least partly on the basis of the depth of their voice. In some species there are satellite males who have no territory and do not call. They may intercept females that are approaching a calling male or take over a vacated territory. Calling is an energy-sapping activity. Sometimes the two roles are reversed and a calling male gives up its territory and becomes a satellite. In explosive breeders, the first male that finds a suitable breeding location, such as a temporary pool, calls loudly and other frogs of both sexes converge on the pool. Explosive breeders tend to call in unison creating a chorus that can be heard from far away. The spadefoot toads (Scaphiopus spp.) of North America fall into this category. Mate selection and courtship is not as important as speed in reproduction. In some years, suitable conditions may not occur and the frogs may go for two or more years without breeding. Some female New Mexico spadefoot toads (Spea multiplicata) only spawn half of the available eggs at a time, perhaps retaining some in case a better reproductive opportunity arises later. At the breeding site, the male mounts the female and grips her tightly round the body. Typically, amplexus takes place in the water, the female releases her eggs and the male covers them with sperm; fertilisation is external. In many species such as the Great Plains toad (Bufo cognatus), the male restrains the eggs with his back feet, holding them in place for about three minutes. Members of the West African genus Nimbaphrynoides are unique among frogs in that they are viviparous; Limnonectes larvaepartus, Eleutherodactylus jasperi and members of the Tanzanian genus Nectophrynoides are the only frogs known to be ovoviviparous. In these species, fertilisation is internal and females give birth to fully developed juvenile frogs, except L. larvaepartus, which give birth to tadpoles. Life cycle Eggs / frogspawn Frogs may lay their in eggs as clumps, surface films, strings, or individually. Around half of species deposit eggs in water, others lay eggs in vegetation, on the ground or in excavations. The tiny yellow-striped pygmy eleuth (Eleutherodactylus limbatus) lays eggs singly, burying them in moist soil. The smoky jungle frog (Leptodactylus pentadactylus) makes a nest of foam in a hollow. The eggs hatch when the nest is flooded, or the tadpoles may complete their development in the foam if flooding does not occur. The red-eyed treefrog (Agalychnis callidryas) deposits its eggs on a leaf above a pool and when they hatch, the larvae fall into the water below. In certain species, such as the wood frog (Rana sylvatica), symbiotic unicellular green algae are present in the gelatinous material. It is thought that these may benefit the developing larvae by providing them with extra oxygen through photosynthesis. The interior of globular egg clusters of the wood frog has also been found to be up to 6 °C (11 °F) warmer than the surrounding water and this speeds up the development of the larvae. The larvae developing in the eggs can detect vibrations caused by nearby predatory wasps or snakes, and will hatch early to avoid being eaten. In general, the length of the egg stage depends on the species and the environmental conditions. Aquatic eggs normally hatch within one week when the capsule splits as a result of enzymes released by the developing larvae. Direct development, where eggs hatch into juveniles like small adults, is also known in many frogs, for example, Ischnocnema henselii, Eleutherodactylus coqui, and Raorchestes ochlandrae and Raorchestes chalazodes. Tadpoles The larvae that emerge from the eggs are known as tadpoles (or occasionally polliwogs). Tadpoles lack eyelids and limbs, and have cartilaginous skeletons, gills for respiration (external gills at first, internal gills later), and tails they use for swimming. As a general rule, free-living larvae are fully aquatic, but at least one species (Nannophrys ceylonensis) has semiterrestrial tadpoles which live among wet rocks. From early in its development, a gill pouch covers the tadpole's gills and front legs. The lungs soon start to develop and are used as an accessory breathing organ. Some species go through metamorphosis while still inside the egg and hatch directly into small frogs. Tadpoles lack true teeth, but the jaws in most species have two elongated, parallel rows of small, keratinized structures called keradonts in their upper jaws. Their lower jaws usually have three rows of keradonts surrounded by a horny beak, but the number of rows can vary and the exact arrangements of mouth parts provide a means for species identification. In the Pipidae, with the exception of Hymenochirus, the tadpoles have paired anterior barbels, which make them resemble small catfish. Their tails are stiffened by a notochord, but does not contain any bony or cartilaginous elements except for a few vertebrae at the base which forms the urostyle during metamorphosis. This has been suggested as an adaptation to their lifestyles; because the transformation into frogs happens very fast, the tail is made of soft tissue only, as bone and cartilage take a much longer time to be broken down and absorbed. The tail fin and tip is fragile and will easily tear, which is seen as an adaptation to escape from predators which try to grasp them by the tail. Tadpoles are typically herbivorous, feeding mostly on algae, including diatoms filtered from the water through the gills. Some species are carnivorous at the tadpole stage, eating insects, smaller tadpoles, and fish. The Cuban tree frog (Osteopilus septentrionalis) is one of a number of species in which the tadpoles can be cannibalistic. Tadpoles that develop legs early may be eaten by the others, so late developers may have better long-term survival prospects. Tadpoles are highly vulnerable to being eaten by fish, newts, predatory diving beetles, and birds, particularly water birds, such as storks and herons and domestic ducks. Some tadpoles, including those of the cane toad (Rhinella marina), are poisonous. The tadpole stage may be as short as a week in explosive breeders or it may last through one or more winters followed by metamorphosis in the spring. Metamorphosis At the end of the tadpole stage, a frog undergoes metamorphosis in which its body makes a sudden transition into the adult form. This metamorphosis typically lasts only 24 hours, and is initiated by production of the hormone thyroxine. This causes different tissues to develop in different ways. The principal changes that take place include the development of the lungs and the disappearance of the gills and gill pouch, making the front legs visible. The lower jaw transforms into the big mandible of the carnivorous adult, and the long, spiral gut of the herbivorous tadpole is replaced by the typical short gut of a predator. Homeostatic feedback control of food intake is largely absent, making tadpoles eat constantly when food is present. But shortly before and during metamorphosis the sensation of hunger is suppressed, and they stop eating while their gut and internal organs are reorganised and prepared for a different diet. Also the gut microbiota changes, from being similar to that of fish to resembling that of amniotes. Exceptions are carnivorous tadpoles like Lepidobatrachus laevis, which has a gut already adapted to a diet similar to that of adults. These continue to eat during metamorphosis. The nervous system becomes adapted for hearing and stereoscopic vision, and for new methods of locomotion and feeding. The eyes are repositioned higher up on the head and the eyelids and associated glands are formed. The eardrum, middle ear, and inner ear are developed. The skin becomes thicker and tougher, the lateral line system is lost, and skin glands are developed. The final stage is the disappearance of the tail, but this takes place rather later, the tissue being used to produce a spurt of growth in the limbs. Frogs are at their most vulnerable to predators when they are undergoing metamorphosis. At this time, the tail is being lost and locomotion by means of limbs is only just becoming established. Adults Adult frogs may live in or near water, but few are fully aquatic. Almost all frog species are carnivorous as adults, preying on invertebrates, including insects, crabs, spiders, mites, worms, snails, and slugs. A few of the larger ones may eat other frogs, small mammals and reptiles, and fish. A few species also eat plant matter; the tree frog Xenohyla truncata is partly herbivorous, its diet including a large proportion of fruit, floral structures and nectar. Leptodactylus mystaceus has been found to eat plants, and folivory occurs in Euphlyctis hexadactylus, with plants constituting 79.5% of its diet by volume. Many frogs use their sticky tongues to catch prey, while others simply grab them with their mouths. Adult frogs are themselves attacked by many predators. The northern leopard frog (Rana pipiens) is eaten by herons, hawks, fish, large salamanders, snakes, raccoons, skunks, mink, bullfrogs, and other animals. Frogs are primary predators and an important part of the food web. Being cold-blooded, they make efficient use of the food they eat with little energy being used for metabolic processes, while the rest is transformed into biomass. They are themselves eaten by secondary predators and are the primary terrestrial consumers of invertebrates, most of which feed on plants. By reducing herbivory, they play a part in increasing the growth of plants and are thus part of a delicately balanced ecosystem. Little is known about the longevity of frogs and toads in the wild, but some can live for many years. Skeletochronology is a method of examining bones to determine age. Using this method, the ages of mountain yellow-legged frogs (Rana muscosa) were studied, the phalanges of the toes showing seasonal lines where growth slows in winter. The oldest frogs had ten bands, so their age was believed to be 14 years, including the four-year tadpole stage. Captive frogs and toads have been recorded as living for up to 40 years, an age achieved by a European common toad (Bufo bufo). The cane toad (Rhinella marina) has been known to survive 24 years in captivity, and the American bullfrog (Rana catesbeiana) 14 years. Frogs from temperate climates hibernate during the winter, and four species are known to be able to withstand freezing during this time, including the wood frog (Rana sylvatica). Parental care Although care of offspring is poorly understood in frogs, up to an estimated 20% of amphibian species may care for their young in some way. The evolution of parental care in frogs is driven primarily by the size of the water body in which they breed. Those that breed in smaller water bodies tend to have greater and more complex parental care behaviour. Because predation of eggs and larvae is high in large water bodies, some frog species started to lay their eggs on land. Once this happened, the desiccating terrestrial environment demands that one or both parents keep them moist to ensure their survival. The subsequent need to transport hatched tadpoles to a water body required an even more intense form of parental care. In small pools, predators are mostly absent and competition between tadpoles becomes the variable that constrains their survival. Certain frog species avoid this competition by making use of smaller phytotelmata (water-filled leaf axils or small woody cavities) as sites for depositing a few tadpoles. While these smaller rearing sites are free from competition, they also lack sufficient nutrients to support a tadpole without parental assistance. Frog species that changed from the use of larger to smaller phytotelmata have evolved a strategy of providing their offspring with nutritive but unfertilised eggs. The female strawberry poison-dart frog (Oophaga pumilio) lays her eggs on the forest floor. The male frog guards them from predation and carries water in his cloaca to keep them moist. When they hatch, the female moves the tadpoles on her back to a water-holding bromeliad or other similar water body, depositing just one in each location. She visits them regularly and feeds them by laying one or two unfertilised eggs in the phytotelma, continuing to do this until the young are large enough to undergo metamorphosis. The granular poison frog (Oophaga granulifera) looks after its tadpoles in a similar way. Many other diverse forms of parental care are seen in frogs. The tiny male Colostethus subpunctatus stands guard over his egg cluster, laid under a stone or log. When the eggs hatch, he transports the tadpoles on his back to a temporary pool, where he partially immerses himself in the water and one or more tadpoles drop off. He then moves on to another pool. The male common midwife toad (Alytes obstetricans) carries the eggs around with him attached to his hind legs. He keeps them damp in dry weather by immersing himself in a pond, and prevents them from getting too wet in soggy vegetation by raising his hindquarters. After three to six weeks, he travels to a pond and the eggs hatch into tadpoles. The tungara frog (Physalaemus pustulosus) builds a floating nest from foam to protect its eggs from predation. The foam is made from proteins and lectins, and seems to have antimicrobial properties. Several pairs of frogs may form a colonial nest on a previously built raft. The eggs are laid in the centre, followed by alternate layers of foam and eggs, finishing with a foam capping. Some frogs protect their offspring inside their own bodies. Both male and female pouched frogs (Assa darlingtoni) guard their eggs, which are laid on the ground. When the eggs hatch, the male lubricates his body with the jelly surrounding them and immerses himself in the egg mass. The tadpoles wriggle into skin pouches on his side, where they develop until they metamorphose into juvenile frogs. The female gastric-brooding frog (Rheobatrachus sp.) from Australia, now probably extinct, swallows her fertilised eggs, which then develop inside her stomach. She ceases to feed and stops secreting stomach acid. The tadpoles rely on the yolks of the eggs for nourishment. After six or seven weeks, they are ready for metamorphosis. The mother regurgitates the tiny frogs, which hop away from her mouth. The female Darwin's frog (Rhinoderma darwinii) from Chile lays up to 40 eggs on the ground, where they are guarded by the male. When the tadpoles are about to hatch, they are engulfed by the male, which carries them around inside his much-enlarged vocal sac. Here they are immersed in a frothy, viscous liquid that contains some nourishment to supplement what they obtain from the yolks of the eggs. They remain in the sac for seven to ten weeks before undergoing metamorphosis, after which they move into the male's mouth and emerge. Defence At first sight, frogs seem rather defenceless because of their small size, slow movement, thin skin, and lack of defensive structures, such as spines, claws or teeth. Many use camouflage to avoid detection, the skin often being spotted or streaked in neutral colours that allow a stationary frog to merge into its surroundings. Some can make prodigious leaps, often into water, that help them to evade potential attackers, while many have other defensive adaptations and strategies. The skin of many frogs contains mild toxic substances called bufotoxins to make them unpalatable to potential predators. Most toads and some frogs have large poison glands, the parotoid glands, located on the sides of their heads behind the eyes and other glands elsewhere on their bodies. These glands secrete mucus and a range of toxins that make frogs slippery to hold and distasteful or poisonous. If the noxious effect is immediate, the predator may cease its action and the frog may escape. If the effect develops more slowly, the predator may learn to avoid that species in future. Poisonous frogs tend to advertise their toxicity with bright colours, an adaptive strategy known as aposematism. The poison dart frogs in the family Dendrobatidae do this. They are typically red, orange, or yellow, often with contrasting black markings on their bodies. Allobates zaparo is not poisonous, but mimics the appearance of two different toxic species with which it shares a common range in an effort to deceive predators. Other species, such as the European fire-bellied toad (Bombina bombina), have their warning colour underneath. They "flash" this when attacked, adopting a pose that exposes the vivid colouring on their bellies. Some frogs, such as the poison dart frogs, are especially toxic. The native peoples of South America extract poison from these frogs to apply to their weapons for hunting, although few species are toxic enough to be used for this purpose. At least two non-poisonous frog species in tropical America (Eleutherodactylus gaigei and Lithodytes lineatus) mimic the colouration of dart poison frogs for self-protection. Some frogs obtain poisons from the ants and other arthropods they eat. Others, such as the Australian corroboree frogs (Pseudophryne corroboree and Pseudophryne pengilleyi), can synthesize the alkaloids themselves. The chemicals involved may be irritants, hallucinogens, convulsants, nerve poisons or vasoconstrictors. Many predators of frogs have become adapted to tolerate high levels of these poisons, but other creatures, including humans who handle the frogs, may be severely affected. Some frogs use bluff or deception. The European common toad (Bufo bufo) adopts a characteristic stance when attacked, inflating its body and standing with its hindquarters raised and its head lowered. The bullfrog (Rana catesbeiana) crouches down with eyes closed and head tipped forward when threatened. This places the parotoid glands in the most effective position, the other glands on its back begin to ooze noxious secretions and the most vulnerable parts of its body are protected. Another tactic used by some frogs is to "scream", the sudden loud noise tending to startle the predator. The grey tree frog (Hyla versicolor) makes an explosive sound that sometimes repels the shrew Blarina brevicauda. Although toads are avoided by many predators, the common garter snake (Thamnophis sirtalis) regularly feeds on them. The strategy employed by juvenile American toads (Bufo americanus) on being approached by a snake is to crouch down and remain immobile. This is usually successful, with the snake passing by and the toad remaining undetected. If it is encountered by the snake's head, however, the toad hops away before crouching defensively. Distribution Frogs live on every continent except Antarctica, but they are not present on certain islands, especially those far away from continental land masses. Many species are isolated in restricted ranges by changes of climate or inhospitable territory, such as stretches of sea, mountain ridges, deserts, forest clearance, road construction, or other human-made barriers. Usually, a greater diversity of frogs occurs in tropical areas than in temperate regions, such as Europe. Some frogs inhabit arid areas, such as deserts, and rely on specific adaptations to survive. Members of the Australian genus Cyclorana bury themselves underground where they create a water-impervious cocoon in which to aestivate during dry periods. Once it rains, they emerge, find a temporary pool, and breed. Egg and tadpole development is very fast compared with those of most other frogs, so breeding can be completed before the pond dries up. Some frog species are adapted to a cold environment. The wood frog (Rana sylvatica), whose habitat extends into the Arctic Circle, buries itself in the ground during winter. Although much of its body freezes during this time, it maintains a high concentration of glucose in its vital organs, which protects them from damage. Conservation In 2006, of 4,035 species of amphibians that depend on water during some lifecycle stage, 1,356 (33.6%) were considered to be threatened. This is likely to be an underestimate because it excludes 1,427 species for which evidence was insufficient to assess their status. Frog populations have declined dramatically since the 1950s. More than one-third of frog species are considered to be threatened with extinction, and more than 120 species are believed to have become extinct since the 1980s. Among these species are the gastric-brooding frogs of Australia and the golden toad of Costa Rica. The latter is of particular concern to scientists because it inhabited the pristine Monteverde Cloud Forest Reserve and its population crashed in 1987, along with about 20 other frog species in the area. This could not be linked directly to human activities, such as deforestation, and was outside the range of normal fluctuations in population size. Elsewhere, habitat loss is a significant cause of frog population decline, as are pollutants, climate change, increased UVB radiation, and the introduction of non-native predators and competitors. A Canadian study conducted in 2006 suggested heavy traffic in their environment was a larger threat to frog populations than was habitat loss. Emerging infectious diseases, including chytridiomycosis and ranavirus, are also devastating populations. Many environmental scientists believe amphibians, including frogs, are good biological indicators of broader ecosystem health because of their intermediate positions in food chains, their permeable skins, and typically biphasic lives (aquatic larvae and terrestrial adults). It appears that species with both aquatic eggs and larvae are most affected by the decline, while those with direct development are the most resistant. Frog mutations and genetic defects have increased since the 1990s. These often include missing legs or extra legs. Various causes have been identified or hypothesized, including an increase in ultraviolet radiation affecting the spawn on the surface of ponds, chemical contamination from pesticides and fertilizers, and parasites such as the trematode Ribeiroia ondatrae. Probably all these are involved in a complex way as stressors, environmental factors contributing to rates of disease, and vulnerability to attack by parasites. Malformations impair mobility and the individuals may not survive to adulthood. An increase in the number of frogs eaten by birds may actually increase the likelihood of parasitism of other frogs, because the trematode's complex lifecycle includes the ramshorn snail and several intermediate hosts such as birds. In a few cases, captive breeding programs have been established and have largely been successful. The World Association of Zoos and Aquariums named 2008 as the "Year of the Frog" in order to draw attention to the conservation issues faced by them. The cane toad (Rhinella marina) is a very adaptable species native to South and Central America. In the 1930s, it was introduced into Puerto Rico, and later various other islands in the Pacific and Caribbean region, as a biological pest control agent. In 1935, 3000 toads were liberated in the sugar cane fields of Queensland, Australia, in an attempt to control cane beetles such as Dermolepida albohirtum, the larvae of which damage and kill the canes. Initial results in many of these countries were positive, but it later became apparent that the toads upset the ecological balance in their new environments. They bred freely, competed with native frog species, ate bees and other harmless native invertebrates, had few predators in their adopted habitats, and poisoned pets, carnivorous birds, and mammals. In many of these countries, they are now regarded both as pests and invasive species, and scientists are looking for a biological method to control them. Human uses Culinary Frog legs are eaten by humans in many parts of the world. Indonesia is the world's largest exporter of frog meat, exporting more than 5,000 tonnes of frog meat each year, mostly to France, Belgium and Luxembourg. Originally, they were supplied from local wild populations, but overexploitation led to a diminution in the supply. This resulted in the development of frog farming and a global trade in frogs. The main importing countries are France, Belgium, Luxembourg, and the United States, while the chief exporting nations are Indonesia and China. The annual global trade in the American bullfrog (Rana catesbeiana), mostly farmed in China, varies between 1200 and 2400 tonnes. The mountain chicken frog, so-called as it tastes of chicken, is now endangered, in part due to human consumption, and was a major food choice of the Dominicans. Raccoon, opossum, partridges, prairie chicken, and frogs were among the fare Mark Twain recorded as part of American cuisine. Scientific research In November 1970, NASA sent two bullfrogs into space for six days during the Orbiting Frog Otolith mission to test weightlessness. Frogs are used for dissections in high school and university anatomy classes, often first being injected with coloured substances to enhance contrasts among the biological systems. This practice is declining due to animal welfare concerns, and "digital frogs" are now available for virtual dissection. Frogs have served as experimental animals throughout the history of science. Eighteenth-century biologist Luigi Galvani discovered the link between electricity and the nervous system by studying frogs. He created one of the first tools for measuring electric current out of a frog leg. In 1852, H. F. Stannius used a frog's heart in a procedure called a Stannius ligature to demonstrate the ventricle and atria beat independently of each other and at different rates. The African clawed frog or platanna (Xenopus laevis) was first widely used in laboratories in pregnancy tests in the first half of the 20th century. A sample of urine from a pregnant woman injected into a female frog induces it to lay eggs, a discovery made by English zoologist Lancelot Hogben. This is because a hormone, human chorionic gonadotropin, is present in substantial quantities in the urine of women during pregnancy. In 1952, Robert Briggs and Thomas J. King cloned a frog by somatic cell nuclear transfer. This same technique was later used to create Dolly the sheep, and their experiment was the first time a successful nuclear transplantation had been accomplished in higher animals. Frogs are used in cloning research and other branches of embryology. Although alternative pregnancy tests have been developed, biologists continue to use Xenopus as a model organism in developmental biology because their embryos are large and easy to manipulate, they are readily obtainable, and can easily be kept in the laboratory. Xenopus laevis is increasingly being displaced by its smaller relative, Xenopus tropicalis, which reaches its reproductive age in five months rather than the one to two years for X. laevis, thus facilitating faster studies across generations. Genomes of Xenopus laevis, X. tropicalis, Rana catesbeiana, Rhinella marina, and Nanorana parkeri have been sequenced and deposited in the NCBI Genome database. Pharmaceutical Because frog toxins are extraordinarily diverse, they have raised the interest of biochemists as a "natural pharmacy". The alkaloid epibatidine, a painkiller 200 times more potent than morphine, is made by some species of poison dart frogs. Other chemicals isolated from the skins of frogs may offer resistance to HIV infection. Dart poisons are under active investigation for their potential as therapeutic drugs. It has long been suspected that pre-Columbian Mesoamericans used a toxic secretion produced by the cane toad as a hallucinogen, but more likely they used substances secreted by the Colorado River toad (Bufo alvarius). These contain bufotenin (5-MeO-DMT), a psychoactive compound that has been used in modern times as a recreational drug. Typically, the skin secretions are dried and then smoked. Illicit drug use by licking the skin of a toad has been reported in the media, but this may be an urban myth. Exudations from the skin of the golden poison frog (Phyllobates terribilis) are traditionally used by native Colombians to poison the darts they use for hunting. The tip of the projectile is rubbed over the back of the frog and the dart is launched from a blowgun. The combination of the two alkaloid toxins batrachotoxin and homobatrachotoxin is so powerful, one frog contains enough poison to kill an estimated 22,000 mice. Two other species, the Kokoe poison dart frog (Phyllobates aurotaenia) and the black-legged dart frog (Phyllobates bicolor) are also used for this purpose. These are less toxic and less abundant than the golden poison frog. They are impaled on pointed sticks and may be heated over a fire to maximise the quantity of poison that can be transferred to the dart. Cultural significance Frogs have been featured in mythology, fairy tales and popular culture. In traditional Chinese myths, the world rests on a giant frog, who would try to swallow the moon, causing the lunar eclipse. Frogs have been featured in religion, folklore, and popular culture. The ancient Egyptians depicted the god Heqet, protector of newborns, with the head of a frog. For the Mayans, frogs represented water, crops, fertility and birth and were associated with the god Chaac. In the Bible, Moses unleashes a plague of frogs on the Egyptians. Medieval Europeans associated frogs and toads with evil and witchcraft. The Brothers Grimm fairy tale The Frog Prince features a princess taking in a frog and it turning into a handsome prince. In modern culture, frogs may take a comedic or hapless role, such as Mr. Toad of the 1908 novel The Wind in the Willows, Michigan J. Frog of Warner Bros. Cartoons, the Muppet Kermit the Frog and in the game Frogger.	Biology and health sciences	Amphibians	null
38499	https://en.wikipedia.org/wiki/Toad	Toad	Toad is a common name for certain frogs, especially of the family Bufonidae, that are characterized by dry, leathery skin, short legs, and large bumps covering the parotoid glands. In popular culture (folk taxonomy), toads are associated with drier, rougher skin and more terrestrial habitats. However, this distinction does not align precisely with scientific taxonomy. List of toad families In scientific taxonomy, toads include the true toads (Bufonidae) and various other terrestrial or warty-skinned frogs. Non-bufonid "toads" can be found in the families: Bombinatoridae (fire-bellied toads and jungle toads) Calyptocephalellidae (helmeted water toad and false toads) Discoglossidae (midwife toads) Myobatrachidae (Australian toadlets) Pelobatidae (European spadefoot toad) Rhinophrynidae (burrowing toads) Scaphiopodidae (American spadefoot toads) Microhylidae (narrowmouth toads) Biology Usually the largest of the bumps on a toad's skin are those that cover the parotoid glands. The bumps are commonly called warts, but they have nothing to do with pathologic warts, being fixed in size, present on healthy specimens, and not caused by infection. It is a myth that handling toads causes warts. Toads travel from non-breeding to breeding areas of ponds and lakes. Bogert (1947) suggests that the toads' call is the most important cue in the homing to ponds. Toads, like many amphibians, exhibit breeding site fidelity (philopatry). Individual American toads return to their natal ponds to breed, making it likely they will encounter siblings when seeking potential mates. Although inbred examples within a species are possible, siblings rarely mate. Toads recognize and avoid mating with close kin. Advertisement vocalizations given by males appear to serve as cues by which females recognize kin. Kin recognition thus allows avoidance of inbreeding and consequent inbreeding depression. Habitat In the United Kingdom, common toads often climb trees to hide in hollows or in nest boxes. Cultural depictions In Kenneth Grahame's novel The Wind in the Willows (1908), Mr. Toad is a likeable and popular, if selfish and narcissistic, comic character. Mr. Toad reappears as the lead character in A. A. Milne's play Toad of Toad Hall (1929), based on the book. In Chinese culture, the Money Toad (or Frog) Jin Chan appears as a feng shui charm for prosperity.	Biology and health sciences	Frogs and toads	Animals
38546	https://en.wikipedia.org/wiki/Nutmeg	Nutmeg	Nutmeg is the seed, or the ground spice derived from that seed, of several tree species of the genus Myristica; fragrant nutmeg or true nutmeg (M. fragrans) is a dark-leaved evergreen tree cultivated for two spices derived from its fruit: nutmeg, from its seed, and mace, from the seed covering. It is also a commercial source of nutmeg essential oil and nutmeg butter. Maluku's Banda Islands are the main producer of nutmeg and mace, and the true nutmeg tree is native to the islands. If consumed in amounts exceeding its typical use as a spice, nutmeg powder may produce allergic reactions, cause contact dermatitis, or have psychoactive effects. Although used in traditional medicine for treating various disorders, nutmeg has no scientifically confirmed medicinal value. Conifers of the genus Torreya, commonly known as the nutmeg yews, have edible seeds of similar appearance, but are not closely related to M. fragrans, and are not used as a spice. Common nutmeg Nutmeg is the spice made by grinding the seed of the fragrant nutmeg tree (Myristica fragrans) into powder. The spice has a distinctive pungent fragrance and a warm, slightly sweet taste; it is used to flavor many kinds of baked goods, confections, puddings, potatoes, meats, sausages, sauces, vegetables, and such beverages as eggnog. The seeds are dried gradually in the sun over a period of 15 to 30 weeks. During this time, the nutmeg shrinks away from its hard seed coat until the kernels rattle in their shells when shaken. The shell is then broken with a wooden club and the nutmegs are picked out. Dried nutmegs are greenish brown ovoids with furrowed surfaces. The nutmegs are roughly egg-shaped, about long and wide, weighing dried. Two other species of genus Myristica with different flavors, M. malabarica and M. argentea, are sometimes used to adulterate nutmeg as a spice. Mace Mace is the spice made from the reddish seed covering (aril) of the nutmeg seed. Its flavour is similar to that of nutmeg but more delicate; it is used to flavour baked goods, meat, fish, and vegetables, and in preserving and pickling. In the processing of mace, the crimson-colored aril is removed from the nutmeg seed that it envelops and is flattened out and dried for 10 to 14 days. Its color changes to pale yellow, orange, or tan. Whole dry mace consists of flat pieces—smooth, horn-like, and brittle—about long. Botany and cultivation The most important commercial species is the common, true or fragrant nutmeg, M. fragrans (Myristicaceae), native to the Moluccas (or Spice Islands) of Indonesia. It is also cultivated on Penang Island in Malaysia, in the Caribbean, especially in Grenada, and in Kerala, a state formerly known as Malabar in ancient writings as the hub of spice trading, in southern India. In the 17th-century work Hortus Botanicus Malabaricus, Hendrik van Rheede records that Indians learned the usage of nutmeg from the Indonesians through ancient trade routes. Nutmeg trees are dioecious plants (individual plants are either male or female), which are propagated sexually from seeds and asexually from cuttings or grafting. Sexual propagation yields 50% male seedlings, which are unproductive. Because no reliable method has been found for determining plant sex before flowering in the sixth to eighth year, and sexual reproduction bears inconsistent yields, grafting is the preferred method of propagation. Epicotyl grafting (a variation of cleft grafting using seedlings), approach grafting, and patch budding have proved successful, with epicotyl grafting being the most widely adopted standard. Air layering is an alternative though not preferred method because of its low (35–40%) success rate. The first harvest of nutmeg trees takes place 7–9 years after planting, and the trees reach full production after 20 years. In the Banda Islands where the nutmeg is endemic, there is a symbiotic relationship between the Kenari nut tree (Canarium indicum) and the nutmeg (Myristica fragrans), the former providing the nutmeg with shade and serving as a wind-break from the strong winds. Uses Spice Nutmeg and mace have similar sensory qualities, with nutmeg having a slightly sweeter and mace a more delicate flavour. Mace is often preferred in light dishes for the bright orange, saffron-like hue it imparts. Nutmeg is used for flavouring many dishes. Whole nutmeg can also be ground at home using a grater specifically designed for nutmeg or a multi-purpose grating tool. In Indonesian cuisine, nutmeg is used in dishes, such as spicy soups including variants of soto, konro, oxtail soup, sup iga (ribs soup), bakso, and sup kambing. It is also used in gravy for meat dishes, such as semur, beef stew, ribs with tomato, and European derived dishes such as bistik (beef steak), rolade (minced meat roll), and bistik lidah (beef tongue steak). In Indian cuisine, nutmeg is used in many sweet, as well as savoury, dishes. In Kerala Malabar region, grated nutmeg is used in meat preparations and also sparingly added to desserts for the flavour. It may also be used in small quantities in garam masala. In traditional European cuisine, nutmeg and mace are used especially in potato and spinach dishes and in processed meat products; they are also used in soups, sauces, and baked goods. It is also commonly used in rice pudding. In Dutch cuisine, nutmeg is added to vegetables such as Brussels sprouts, cauliflower, and string beans. Nutmeg is a traditional ingredient in mulled cider, mulled wine, junket and eggnog. In Scotland, mace and nutmeg are usually both ingredients in haggis. In Italian cuisine, nutmeg is used as part of the stuffing for many regional meat-filled dumplings like tortellini, as well as for the traditional meatloaf. Nutmeg is a common spice for pumpkin pie and in recipes for other winter squashes, such as baked acorn squash. In the Caribbean, nutmeg is often used in drinks, such as the Bushwacker, Painkiller, and Barbados rum punch. Typically, it is a sprinkle on top of the drink. Fruit The pericarp (fruit covering) is used to make jam, or is finely sliced, cooked with sugar, and crystallised to make a fragrant candy. Sliced nutmeg fruit flesh is made as manisan (sweets), either wet, which is seasoned in sugary syrup liquid, or dry coated with sugar, a dessert called manisan pala in Indonesia. In Penang cuisine, dried, shredded nutmeg rind with sugar coating is used as toppings on the uniquely Penang ais kacang. The flesh of the nutmeg fruit is also blended, in the fresh state, into a type of smoothie (white in colour and having a fresh, ‘green’, tangy taste); or boiled, resulting in a brown liquid, much sweeter in taste, which is used in the preparation of iced drinks. In Kerala Malabar region of India, it is used for juice, pickles and chutney. Essential oil The essential oil obtained by steam distillation of ground nutmeg is used in the perfumery and pharmaceutical industries. The volatile fraction contains dozens of terpenes and phenylpropanoids, including -pinene, limonene, -borneol, -terpineol, geraniol, safrole, and myristicin. In its pure form, myristicin is a toxin, and consumption of excessive amounts of nutmeg can result in myristicin poisoning. The oil is colorless or light yellow, and smells and tastes of nutmeg. It is used as a natural food flavoring in baked goods, syrups, beverages, and sweets. It is used to replace ground nutmeg, as it leaves no particles in the food. The essential oil is also used in the manufacturing of toothpaste and cough syrups. Nutmeg butter Nutmeg butter is obtained from the nut by expression. It is semisolid, reddish-brown in colour, and has the taste and smell of nutmeg itself. About 75% (by weight) of nutmeg butter is trimyristin, which can be turned into myristic acid, a 14-carbon fatty acid, which can be used as a replacement for cocoa butter, can be mixed with other fats like cottonseed oil or palm oil, and has applications as an industrial lubricant. History The earliest evidence of use of nutmeg comes in the form of 3,500-year-old potsherd residues from the island of Pulau Ai, one of the Banda Islands in eastern Indonesia. The Banda Islands consist of eleven small volcanic islands, and are part of the larger Maluku Islands group. These islands were the only source of nutmeg and mace production until the mid-19th century. It was one of the spices traded over the Austronesian maritime spice trade network since at least 1500 BCE. In the sixth century AD, nutmeg use spread to India, then further west to Constantinople. By the 13th century, Arab traders had pinpointed the origin of nutmeg to the Banda Islands, but kept this location a secret from European traders. Colonial era The Banda Islands became the scene of the earliest European ventures in Asia, to get a grip on the spice trade. In August 1511, Afonso de Albuquerque conquered Malacca, which at the time was the hub of Asian trade, on behalf of the king of Portugal. In November of the same year, after having secured Malacca and learning of Banda's location, Albuquerque sent an expedition of three ships led by his friend António de Abreu to find it. Malay pilots guided them via Java, the Lesser Sundas, and Ambon to the Banda Islands, arriving in early 1512. The first Europeans to reach the Banda Islands, the expedition remained for about a month, buying and filling their ships with Banda's nutmeg and mace, and with cloves in which Banda had a thriving entrepôt trade. An early account of Banda is in Suma Oriental, a book written by the Portuguese apothecary Tomé Pires, based in Malacca from 1512 to 1515. Full control of this trade by the Portuguese was not possible, and they remained participants without a foothold in the islands. In order to obtain a monopoly on the production and trade of nutmeg, the Dutch East India Company (VOC) waged a bloody battle with the Bandanese in 1621. Historian Willard Hanna estimated that before this struggle the islands were populated by approximately 15,000 people, and only 1,000 were left (the Bandanese were killed, starved while fleeing, exiled, or sold as slaves). The Company constructed a comprehensive nutmeg plantation system on the islands during the 17th century. As a result of the Dutch interregnum during the Napoleonic Wars, the British invaded and temporarily took control of the Banda Islands from the Dutch and transplanted nutmeg trees, complete with soil, to Sri Lanka, Penang, Bencoolen, and Singapore. From these locations they were transplanted to their other colonial holdings elsewhere, notably Zanzibar and Grenada. The national flag of Grenada, adopted in 1974, shows a stylised split-open nutmeg fruit. The Dutch retained control of the Spice Islands until World War II. Connecticut may have received its nickname ("the Nutmeg State", "Nutmegger") from the claim that some unscrupulous Connecticut traders would whittle "nutmeg" out of wood, creating a "wooden nutmeg", a term which later came to mean any type of fraud. This narrative may have to do with the issue that one has to grate to obtain the spice powder, not crack a nutmeg, and this may not have been widely known by some purchasers of the product. Production In 2019, global production of nutmeg was 142,000 tonnes, led by Indonesia, Guatemala, and India, having 38,000 to 43,000 tonnes each and a combined 85% of the world total. Psychoactivity and toxicity Although used as a folk treatment for some ailments, nutmeg has no proven medicinal value. Effects Ingested in small amounts as a spice, nutmeg produces no noticeable physiological or neurological response, but in large doses, both raw nutmeg freshly ground from kernels and nutmeg oil have psychoactive effects. Such effects appear to derive from anticholinergic-like hallucinogenic mechanisms attributed to myristicin and elemicin. Myristicin—a monoamine oxidase inhibitor and psychoactive substance—can cause convulsions, palpitations, nausea, eventual dehydration, and generalized body pain when consumed in large amounts. Nutmeg may interact with anxiolytic drugs, produce allergic reactions, cause contact dermatitis, and evoke acute episodes of psychosis. Varying considerably from person to person, nutmeg intoxication may occur with side effects, such as delirium, anxiety, confusion, headaches, nausea, dizziness, dry mouth, eye irritation, and amnesia. Intoxication takes several hours to reach maximum effect, and may last for several days. Incidents of fatal poisoning from nutmeg and myristicin individually are uncommon. Nutmeg poisonings occur by accidental consumption in children and by intentional recreational use. It is used recreationally with the intention of achieving a low-cost high resembling psychedelics, particularly by adolescents, drug users, college students, and prisoners. Relatively large doses of nutmeg are required to produce effects; a majority of reported nutmeg intoxication cases appear to result from recreational use. Playwright and poet William Shakespeare was alleged to use nutmeg for hallucinogenic purposes as nutmeg extract along with cannabis were found in analysis of fragments of his pipe. Toxicity during pregnancy Nutmeg was once considered an abortifacient, but may be safe during pregnancy if used only in flavoring amounts. If consumed in large amounts, nutmeg could cause premature labor and miscarriage. Nutmeg may also interact with pain relievers such as pethidine, so avoiding it during pregnancy is recommended. Toxicity to pets The scent of nutmeg may attract pets, but it can be poisonous if consumed in excess.	Biology and health sciences	Magnoliids	null
38571	https://en.wikipedia.org/wiki/Adenine	Adenine	Adenine (, ) (symbol A or Ade) is a purine nucleotide base. It is one of the nucleobases in the nucleic acids, DNA and RNA. The shape of adenine is complementary to either thymine in DNA or uracil in RNA. In cells adenine, as an independent molecule, is rare. It is almost always covalently bound to become a part of a larger biomolecule. Adenine has a central role in cellular respiration. It is part of adenosine triphosphate which provides the energy that drives and supports most activities in living cells, such as protein synthesis, chemical synthesis, muscle contraction, and nerve impulse propagation. In respiration it also participates as part of the cofactors nicotinamide adenine dinucleotide, flavin adenine dinucleotide, and Coenzyme A. It is also part of adenosine, adenosine monophosphate, cyclic adenosine monophosphate, adenosine diphosphate, and S-adenosylmethionine. Structure Adenine forms several tautomers, compounds that can be rapidly interconverted and are often considered equivalent. However, in isolated conditions, i.e. in an inert gas matrix and in the gas phase, mainly the 9H-adenine tautomer is found. Biosynthesis Purine metabolism involves the formation of adenine and guanine. Both adenine and guanine are derived from the nucleotide inosine monophosphate (IMP), which in turn is synthesized from a pre-existing ribose phosphate through a complex pathway using atoms from the amino acids glycine, glutamine, and aspartic acid, as well as the coenzyme tetrahydrofolate. Manufacturing method Patented August 20, 1968, the current recognized method of industrial-scale production of adenine is a modified form of the formamide method. This method heats up formamide under 120 degree Celsius conditions within a sealed flask for 5 hours to form adenine. The reaction is heavily increased in quantity by using a phosphorus oxychloride (phosphoryl chloride) or phosphorus pentachloride as an acid catalyst and sunlight or ultraviolet conditions. After the 5 hours have passed and the formamide-phosphorus oxychloride-adenine solution cools down, water is put into the flask containing the formamide and now-formed adenine. The water-formamide-adenine solution is then poured through a filtering column of activated charcoal. The water and formamide molecules, being small molecules, will pass through the charcoal and into the waste flask; the large adenine molecules, however, will attach or "adsorb" to the charcoal due to the van der Waals forces that interact between the adenine and the carbon in the charcoal. Because charcoal has a large surface area, it is able to capture the majority of molecules that pass a certain size (greater than water and formamide) through it. To extract the adenine from the charcoal-adsorbed adenine, ammonia gas dissolved in water (aqua ammonia) is poured onto the activated charcoal-adenine structure to liberate the adenine into the ammonia-water solution. The solution containing water, ammonia, and adenine is then left to air dry, with the adenine losing solubility due to the loss of ammonia gas that previously made the solution basic and capable of dissolving adenine, thus causing it to crystallize into a pure white powder that can be stored. Function Adenine is one of the two purine nucleobases (the other being guanine) used in forming nucleotides of the nucleic acids. In DNA, adenine binds to thymine via two hydrogen bonds to assist in stabilizing the nucleic acid structures. In RNA, which is used for protein synthesis, adenine binds to uracil. Adenine forms adenosine, a nucleoside, when attached to ribose, and deoxyadenosine when attached to deoxyribose. It forms adenosine triphosphate (ATP), a nucleoside triphosphate, when three phosphate groups are added to adenosine. Adenosine triphosphate is used in cellular metabolism as one of the basic methods of transferring chemical energy between chemical reactions. ATP is thus a derivative of adenine, adenosine, cyclic adenosine monophosphate, and adenosine diphosphate. {\| class="wikitable left" style="text-align:center" \|- \| \| \|- \| Adenosine, A \| Deoxyadenosine, dA \|} History In older literature, adenine was sometimes called Vitamin B4. Due to it being synthesized by the body and not essential to be obtained by diet, it does not meet the definition of vitamin and is no longer part of the Vitamin B complex. However, two B vitamins, niacin and riboflavin, bind with adenine to form the essential cofactors nicotinamide adenine dinucleotide (NAD) and flavin adenine dinucleotide (FAD), respectively. Hermann Emil Fischer was one of the early scientists to study adenine. It was named in 1885 by Albrecht Kossel after Greek ἀδήν aden "gland", in reference to the pancreas, from which Kossel's sample had been extracted. Experiments performed in 1961 by Joan Oró have shown that a large quantity of adenine can be synthesized from the polymerization of ammonia with five hydrogen cyanide (HCN) molecules in aqueous solution; whether this has implications for the origin of life on Earth is under debate. On August 8, 2011, a report, based on NASA studies with meteorites found on Earth, was published suggesting building blocks of DNA and RNA (adenine, guanine and related organic molecules) may have been formed extraterrestrially in outer space. In 2011, physicists reported that adenine has an "unexpectedly variable range of ionization energies along its reaction pathways" which suggested that "understanding experimental data on how adenine survives exposure to UV light is much more complicated than previously thought"; these findings have implications for spectroscopic measurements of heterocyclic compounds, according to one report.	Biology and health sciences	Nucleic acids	Biology
38579	https://en.wikipedia.org/wiki/Gravity	Gravity	In physics, gravity () is a fundamental interaction primarily observed as mutual attraction between all things that have mass. Gravity is, by far, the weakest of the four fundamental interactions, approximately 1038 times weaker than the strong interaction, 1036 times weaker than the electromagnetic force and 1029 times weaker than the weak interaction. As a result, it has no significant influence at the level of subatomic particles. However, gravity is the most significant interaction between objects at the macroscopic scale, and it determines the motion of planets, stars, galaxies, and even light. On Earth, gravity gives weight to physical objects, and the Moon's gravity is responsible for sublunar tides in the oceans. The corresponding antipodal tide is caused by the inertia of the Earth and Moon orbiting one another. Gravity also has many important biological functions, helping to guide the growth of plants through the process of gravitropism and influencing the circulation of fluids in multicellular organisms. The gravitational attraction between the original gaseous matter in the universe caused it to coalesce and form stars which eventually condensed into galaxies, so gravity is responsible for many of the large-scale structures in the universe. Gravity has an infinite range, although its effects become weaker as objects get farther away. Gravity is most accurately described by the general theory of relativity, proposed by Albert Einstein in 1915, which describes gravity not as a force, but as the curvature of spacetime, caused by the uneven distribution of mass, and causing masses to move along geodesic lines. The most extreme example of this curvature of spacetime is a black hole, from which nothing—not even light—can escape once past the black hole's event horizon. However, for most applications, gravity is well approximated by Newton's law of universal gravitation, which describes gravity as a force causing any two bodies to be attracted toward each other, with magnitude proportional to the product of their masses and inversely proportional to the square of the distance between them. Current models of particle physics imply that the earliest instance of gravity in the universe, possibly in the form of quantum gravity, supergravity or a gravitational singularity, along with ordinary space and time, developed during the Planck epoch (up to 10−43 seconds after the birth of the universe), possibly from a primeval state, such as a false vacuum, quantum vacuum or virtual particle, in a currently unknown manner. Scientists are currently working to develop a theory of gravity consistent with quantum mechanics, a quantum gravity theory, which would allow gravity to be united in a common mathematical framework (a theory of everything) with the other three fundamental interactions of physics. Definitions , also known as gravitational attraction, is the mutual attraction between all masses in the universe. Gravity is the gravitational attraction at the surface of a planet or other celestial body; gravity may also include, in addition to gravitation, the centrifugal force resulting from the planet's rotation . History Ancient world The nature and mechanism of gravity were explored by a wide range of ancient scholars. In Greece, Aristotle believed that objects fell towards the Earth because the Earth was the center of the Universe and attracted all of the mass in the Universe towards it. He also thought that the speed of a falling object should increase with its weight, a conclusion that was later shown to be false. While Aristotle's view was widely accepted throughout Ancient Greece, there were other thinkers such as Plutarch who correctly predicted that the attraction of gravity was not unique to the Earth. Although he did not understand gravity as a force, the ancient Greek philosopher Archimedes discovered the center of gravity of a triangle. He postulated that if two equal weights did not have the same center of gravity, the center of gravity of the two weights together would be in the middle of the line that joins their centers of gravity. Two centuries later, the Roman engineer and architect Vitruvius contended in his De architectura that gravity is not dependent on a substance's weight but rather on its "nature". In the 6th century CE, the Byzantine Alexandrian scholar John Philoponus proposed the theory of impetus, which modifies Aristotle's theory that "continuation of motion depends on continued action of a force" by incorporating a causative force that diminishes over time. In 628 CE, the Indian mathematician and astronomer Brahmagupta proposed the idea that gravity is an attractive force that draws objects to the Earth and used the term gurutvākarṣaṇ to describe it. In the ancient Middle East, gravity was a topic of fierce debate. The Persian intellectual Al-Biruni believed that the force of gravity was not unique to the Earth, and he correctly assumed that other heavenly bodies should exert a gravitational attraction as well. In contrast, Al-Khazini held the same position as Aristotle that all matter in the Universe is attracted to the center of the Earth. Scientific revolution In the mid-16th century, various European scientists experimentally disproved the Aristotelian notion that heavier objects fall at a faster rate. In particular, the Spanish Dominican priest Domingo de Soto wrote in 1551 that bodies in free fall uniformly accelerate. De Soto may have been influenced by earlier experiments conducted by other Dominican priests in Italy, including those by Benedetto Varchi, Francesco Beato, Luca Ghini, and Giovan Bellaso which contradicted Aristotle's teachings on the fall of bodies. The mid-16th century Italian physicist Giambattista Benedetti published papers claiming that, due to specific gravity, objects made of the same material but with different masses would fall at the same speed. With the 1586 Delft tower experiment, the Flemish physicist Simon Stevin observed that two cannonballs of differing sizes and weights fell at the same rate when dropped from a tower. In the late 16th century, Galileo Galilei's careful measurements of balls rolling down inclines allowed him to firmly establish that gravitational acceleration is the same for all objects. Galileo postulated that air resistance is the reason that objects with a low density and high surface area fall more slowly in an atmosphere. In 1604, Galileo correctly hypothesized that the distance of a falling object is proportional to the square of the time elapsed. This was later confirmed by Italian scientists Jesuits Grimaldi and Riccioli between 1640 and 1650. They also calculated the magnitude of the Earth's gravity by measuring the oscillations of a pendulum. Newton's theory of gravitation In 1657, Robert Hooke published his Micrographia, in which he hypothesized that the Moon must have its own gravity. In 1666, he added two further principles: that all bodies move in straight lines until deflected by some force and that the attractive force is stronger for closer bodies. In a communication to the Royal Society in 1666, Hooke wrote Hooke's 1674 Gresham lecture, An Attempt to prove the Annual Motion of the Earth, explained that gravitation applied to "all celestial bodies" In 1684, Newton sent a manuscript to Edmond Halley titled De motu corporum in gyrum ('On the motion of bodies in an orbit'), which provided a physical justification for Kepler's laws of planetary motion. Halley was impressed by the manuscript and urged Newton to expand on it, and a few years later Newton published a groundbreaking book called Philosophiæ Naturalis Principia Mathematica (Mathematical Principles of Natural Philosophy). In this book, Newton described gravitation as a universal force, and claimed that "the forces which keep the planets in their orbs must [be] reciprocally as the squares of their distances from the centers about which they revolve." This statement was later condensed into the following inverse-square law: where is the force, and are the masses of the objects interacting, is the distance between the centers of the masses and is the gravitational constant Newton's Principia was well received by the scientific community, and his law of gravitation quickly spread across the European world. More than a century later, in 1821, his theory of gravitation rose to even greater prominence when it was used to predict the existence of Neptune. In that year, the French astronomer Alexis Bouvard used this theory to create a table modeling the orbit of Uranus, which was shown to differ significantly from the planet's actual trajectory. In order to explain this discrepancy, many astronomers speculated that there might be a large object beyond the orbit of Uranus which was disrupting its orbit. In 1846, the astronomers John Couch Adams and Urbain Le Verrier independently used Newton's law to predict Neptune's location in the night sky, and the planet was discovered there within a day. General relativity Eventually, astronomers noticed an eccentricity in the orbit of the planet Mercury which could not be explained by Newton's theory: the perihelion of the orbit was increasing by about 42.98 arcseconds per century. The most obvious explanation for this discrepancy was an as-yet-undiscovered celestial body, such as a planet orbiting the Sun even closer than Mercury, but all efforts to find such a body turned out to be fruitless. In 1915, Albert Einstein developed a theory of general relativity which was able to accurately model Mercury's orbit. In general relativity, the effects of gravitation are ascribed to spacetime curvature instead of a force. Einstein began to toy with this idea in the form of the equivalence principle, a discovery which he later described as "the happiest thought of my life." In this theory, free fall is considered to be equivalent to inertial motion, meaning that free-falling inertial objects are accelerated relative to non-inertial observers on the ground. In contrast to Newtonian physics, Einstein believed that it was possible for this acceleration to occur without any force being applied to the object. Einstein proposed that spacetime is curved by matter, and that free-falling objects are moving along locally straight paths in curved spacetime. These straight paths are called geodesics. As in Newton's first law of motion, Einstein believed that a force applied to an object would cause it to deviate from a geodesic. For instance, people standing on the surface of the Earth are prevented from following a geodesic path because the mechanical resistance of the Earth exerts an upward force on them. This explains why moving along the geodesics in spacetime is considered inertial. Einstein's description of gravity was quickly accepted by the majority of physicists, as it was able to explain a wide variety of previously baffling experimental results. In the coming years, a wide range of experiments provided additional support for the idea of general relativity. Today, Einstein's theory of relativity is used for all gravitational calculations where absolute precision is desired, although Newton's inverse-square law is accurate enough for virtually all ordinary calculations. Modern research In modern physics, general relativity remains the framework for the understanding of gravity. Physicists continue to work to find solutions to the Einstein field equations that form the basis of general relativity and continue to test the theory, finding excellent agreement in all cases. Einstein field equations The Einstein field equations are a system of 10 partial differential equations which describe how matter affects the curvature of spacetime. The system is often expressed in the form where is the Einstein tensor, is the metric tensor, is the stress–energy tensor, is the cosmological constant, is the Newtonian constant of gravitation and is the speed of light. The constant is referred to as the Einstein gravitational constant. A major area of research is the discovery of exact solutions to the Einstein field equations. Solving these equations amounts to calculating a precise value for the metric tensor (which defines the curvature and geometry of spacetime) under certain physical conditions. There is no formal definition for what constitutes such solutions, but most scientists agree that they should be expressable using elementary functions or linear differential equations. Some of the most notable solutions of the equations include: The Schwarzschild solution, which describes spacetime surrounding a spherically symmetric non-rotating uncharged massive object. For compact enough objects, this solution generated a black hole with a central singularity. At points far away from the central mass, the accelerations predicted by the Schwarzschild solution are practically identical to those predicted by Newton's theory of gravity. The Reissner–Nordström solution, which analyzes a non-rotating spherically symmetric object with charge and was independently discovered by several different researchers between 1916 and 1921. In some cases, this solution can predict the existence of black holes with double event horizons. The Kerr solution, which generalizes the Schwarzchild solution to rotating massive objects. Because of the difficulty of factoring in the effects of rotation into the Einstein field equations, this solution was not discovered until 1963. The Kerr–Newman solution for charged, rotating massive objects. This solution was derived in 1964, using the same technique of complex coordinate transformation that was used for the Kerr solution. The cosmological Friedmann–Lemaître–Robertson–Walker solution, discovered in 1922 by Alexander Friedmann and then confirmed in 1927 by Georges Lemaître. This solution was revolutionary for predicting the expansion of the Universe, which was confirmed seven years later after a series of measurements by Edwin Hubble. It even showed that general relativity was incompatible with a static universe, and Einstein later conceded that he had been wrong to design his field equations to account for a Universe that was not expanding. Today, there remain many important situations in which the Einstein field equations have not been solved. Chief among these is the two-body problem, which concerns the geometry of spacetime around two mutually interacting massive objects, such as the Sun and the Earth, or the two stars in a binary star system. The situation gets even more complicated when considering the interactions of three or more massive bodies (the "n-body problem"), and some scientists suspect that the Einstein field equations will never be solved in this context. However, it is still possible to construct an approximate solution to the field equations in the n-body problem by using the technique of post-Newtonian expansion. In general, the extreme nonlinearity of the Einstein field equations makes it difficult to solve them in all but the most specific cases. Gravity and quantum mechanics Despite its success in predicting the effects of gravity at large scales, general relativity is ultimately incompatible with quantum mechanics. This is because general relativity describes gravity as a smooth, continuous distortion of spacetime, while quantum mechanics holds that all forces arise from the exchange of discrete particles known as quanta. This contradiction is especially vexing to physicists because the other three fundamental forces (strong force, weak force and electromagnetism) were reconciled with a quantum framework decades ago. As a result, modern researchers have begun to search for a theory that could unite both gravity and quantum mechanics under a more general framework. One path is to describe gravity in the framework of quantum field theory, which has been successful to accurately describe the other fundamental interactions. The electromagnetic force arises from an exchange of virtual photons, where the QFT description of gravity is that there is an exchange of virtual gravitons. This description reproduces general relativity in the classical limit. However, this approach fails at short distances of the order of the Planck length, where a more complete theory of quantum gravity (or a new approach to quantum mechanics) is required. Tests of general relativity Testing the predictions of general relativity has historically been difficult, because they are almost identical to the predictions of Newtonian gravity for small energies and masses. Still, since its development, an ongoing series of experimental results have provided support for the theory: In 1919, the British astrophysicist Arthur Eddington was able to confirm the predicted gravitational lensing of light during that year's solar eclipse. Eddington measured starlight deflections twice those predicted by Newtonian corpuscular theory, in accordance with the predictions of general relativity. Although Eddington's analysis was later disputed, this experiment made Einstein famous almost overnight and caused general relativity to become widely accepted in the scientific community. In 1959, American physicists Robert Pound and Glen Rebka performed an experiment in which they used gamma rays to confirm the prediction of gravitational time dilation. By sending the rays down a 74-foot tower and measuring their frequency at the bottom, the scientists confirmed that light is redshifted as it moves towards a source of gravity. The observed redshift also supported the idea that time runs more slowly in the presence of a gravitational field. The time delay of light passing close to a massive object was first identified by Irwin I. Shapiro in 1964 in interplanetary spacecraft signals. In 1971, scientists discovered the first-ever black hole in the galaxy Cygnus. The black hole was detected because it was emitting bursts of x-rays as it consumed a smaller star, and it came to be known as Cygnus X-1. This discovery confirmed yet another prediction of general relativity, because Einstein's equations implied that light could not escape from a sufficiently large and compact object. General relativity states that gravity acts on light and matter equally, meaning that a sufficiently massive object could warp light around it and create a gravitational lens. This phenomenon was first confirmed by observation in 1979 using the 2.1 meter telescope at Kitt Peak National Observatory in Arizona, which saw two mirror images of the same quasar whose light had been bent around the galaxy YGKOW G1. Frame dragging, the idea that a rotating massive object should twist spacetime around it, was confirmed by Gravity Probe B results in 2011. In 2015, the LIGO observatory detected faint gravitational waves, the existence of which had been predicted by general relativity. Scientists believe that the waves emanated from a black hole merger that occurred 1.5 billion light-years away. Specifics Earth's gravity Every planetary body (including the Earth) is surrounded by its own gravitational field, which can be conceptualized with Newtonian physics as exerting an attractive force on all objects. Assuming a spherically symmetrical planet, the strength of this field at any given point above the surface is proportional to the planetary body's mass and inversely proportional to the square of the distance from the center of the body. The strength of the gravitational field is numerically equal to the acceleration of objects under its influence. The rate of acceleration of falling objects near the Earth's surface varies very slightly depending on latitude, surface features such as mountains and ridges, and perhaps unusually high or low sub-surface densities. For purposes of weights and measures, a standard gravity value is defined by the International Bureau of Weights and Measures, under the International System of Units (SI). The force of gravity experienced by objects on Earth's surface is the vector sum of two forces: (a) The gravitational attraction in accordance with Newton's universal law of gravitation, and (b) the centrifugal force, which results from the choice of an earthbound, rotating frame of reference. The force of gravity is weakest at the equator because of the centrifugal force caused by the Earth's rotation and because points on the equator are furthest from the center of the Earth. The force of gravity varies with latitude and increases from about 9.780 m/s2 at the Equator to about 9.832 m/s2 at the poles. Gravitational radiation General relativity predicts that energy can be transported out of a system through gravitational radiation. The first indirect evidence for gravitational radiation was through measurements of the Hulse–Taylor binary in 1973. This system consists of a pulsar and neutron star in orbit around one another. Its orbital period has decreased since its initial discovery due to a loss of energy, which is consistent for the amount of energy loss due to gravitational radiation. This research was awarded the Nobel Prize in Physics in 1993. The first direct evidence for gravitational radiation was measured on 14 September 2015 by the LIGO detectors. The gravitational waves emitted during the collision of two black holes 1.3 billion light years from Earth were measured. This observation confirms the theoretical predictions of Einstein and others that such waves exist. It also opens the way for practical observation and understanding of the nature of gravity and events in the Universe including the Big Bang. Neutron star and black hole formation also create detectable amounts of gravitational radiation. This research was awarded the Nobel Prize in Physics in 2017. Speed of gravity In December 2012, a research team in China announced that it had produced measurements of the phase lag of Earth tides during full and new moons which seem to prove that the speed of gravity is equal to the speed of light. This means that if the Sun suddenly disappeared, the Earth would keep orbiting the vacant point normally for 8 minutes, which is the time light takes to travel that distance. The team's findings were released in Science Bulletin in February 2013. In October 2017, the LIGO and Virgo detectors received gravitational wave signals within 2 seconds of gamma ray satellites and optical telescopes seeing signals from the same direction. This confirmed that the speed of gravitational waves was the same as the speed of light. Anomalies and discrepancies There are some observations that are not adequately accounted for, which may point to the need for better theories of gravity or perhaps be explained in other ways. Extra-fast stars: Stars in galaxies follow a distribution of velocities where stars on the outskirts are moving faster than they should according to the observed distributions of normal matter. Galaxies within galaxy clusters show a similar pattern. Dark matter, which would interact through gravitation but not electromagnetically, would account for the discrepancy. Various modifications to Newtonian dynamics have also been proposed. Accelerated expansion: The expansion of the universe seems to be speeding up. Dark energy has been proposed to explain this. Flyby anomaly: Various spacecraft have experienced greater acceleration than expected during gravity assist maneuvers. The Pioneer anomaly has been shown to be explained by thermal recoil due to the distant sun radiation on one side of the space craft. Alternative theories Historical alternative theories Aristotelian theory of gravity Le Sage's theory of gravitation (1784) also called LeSage gravity but originally proposed by Fatio and further elaborated by Georges-Louis Le Sage, based on a fluid-based explanation where a light gas fills the entire Universe. Ritz's theory of gravitation, Ann. Chem. Phys. 13, 145, (1908) pp. 267–271, Weber–Gauss electrodynamics applied to gravitation. Classical advancement of perihelia. Nordström's theory of gravitation (1912, 1913), an early competitor of general relativity. Kaluza–Klein theory (1921) Whitehead's theory of gravitation (1922), another early competitor of general relativity. Modern alternative theories Brans–Dicke theory of gravity (1961) Induced gravity (1967), a proposal by Andrei Sakharov according to which general relativity might arise from quantum field theories of matter String theory (late 1960s) ƒ(R) gravity (1970) Horndeski theory (1974) Supergravity (1976) In the modified Newtonian dynamics (MOND) (1981), Mordehai Milgrom proposes a modification of Newton's second law of motion for small accelerations The self-creation cosmology theory of gravity (1982) by G.A. Barber in which the Brans–Dicke theory is modified to allow mass creation Loop quantum gravity (1988) by Carlo Rovelli, Lee Smolin, and Abhay Ashtekar Nonsymmetric gravitational theory (NGT) (1994) by John Moffat Tensor–vector–scalar gravity (TeVeS) (2004), a relativistic modification of MOND by Jacob Bekenstein Chameleon theory (2004) by Justin Khoury and Amanda Weltman. Pressuron theory (2013) by Olivier Minazzoli and Aurélien Hees. Conformal gravity Gravity as an entropic force, gravity arising as an emergent phenomenon from the thermodynamic concept of entropy. In the superfluid vacuum theory the gravity and curved spacetime arise as a collective excitation mode of non-relativistic background superfluid. Massive gravity, a theory where gravitons and gravitational waves have a non-zero mass	Physical sciences	Physics	null
8057418	https://en.wikipedia.org/wiki/Quantum%20potential	Quantum potential	The quantum potential or quantum potentiality is a central concept of the de Broglie–Bohm formulation of quantum mechanics, introduced by David Bohm in 1952. Initially presented under the name quantum-mechanical potential, subsequently quantum potential, it was later elaborated upon by Bohm and Basil Hiley in its interpretation as an information potential which acts on a quantum particle. It is also referred to as quantum potential energy, Bohm potential, quantum Bohm potential or Bohm quantum potential. In the framework of the de Broglie–Bohm theory, the quantum potential is a term within the Schrödinger equation which acts to guide the movement of quantum particles. The quantum potential approach introduced by Bohm provides a physically less fundamental exposition of the idea presented by Louis de Broglie: de Broglie had postulated in 1925 that the relativistic wave function defined on spacetime represents a pilot wave which guides a quantum particle, represented as an oscillating peak in the wave field, but he had subsequently abandoned his approach because he was unable to derive the guidance equation for the particle from a non-linear wave equation. The seminal articles of Bohm in 1952 introduced the quantum potential and included answers to the objections which had been raised against the pilot wave theory. The Bohm quantum potential is closely linked with the results of other approaches, in particular relating to works of Erwin Madelung in 1927 and Carl Friedrich von Weizsäcker in 1935. Building on the interpretation of the quantum theory introduced by Bohm in 1952, David Bohm and Basil Hiley in 1975 presented how the concept of a quantum potential leads to the notion of an "unbroken wholeness of the entire universe", proposing that the fundamental new quality introduced by quantum physics is nonlocality. Relation to the Schrödinger equation The Schrödinger equation is re-written using the polar form for the wave function with real-valued functions and , where is the amplitude (absolute value) of the wave function , and its phase. This yields two equations: from the imaginary and real part of the Schrödinger equation follow the continuity equation and the quantum Hamilton–Jacobi equation respectively. Continuity equation The imaginary part of the Schrödinger equation in polar form yields which, provided , can be interpreted as the continuity equation for the probability density and the velocity field Quantum Hamilton–Jacobi equation The real part of the Schrödinger equation in polar form yields a modified Hamilton–Jacobi equation also referred to as quantum Hamilton–Jacobi equation. It differs from the classical Hamilton–Jacobi equation only by the term This term , called quantum potential, thus depends on the curvature of the amplitude of the wave function. In the limit , the function is a solution of the (classical) Hamilton–Jacobi equation; therefore, the function is also called the Hamilton–Jacobi function, or action, extended to quantum physics. Properties Hiley emphasised several aspects that regard the quantum potential of a quantum particle: it is derived mathematically from the real part of the Schrödinger equation under polar decomposition of the wave function, is not derived from a Hamiltonian or other external source, and could be said to be involved in a self-organising process involving a basic underlying field; it does not change if is multiplied by a constant, as this term is also present in the denominator, so that is independent of the magnitude of and thus of field intensity; therefore, the quantum potential fulfils a precondition for nonlocality: it need not fall off as distance increases; it carries information about the whole experimental arrangement in which the particle finds itself. In 1979, Hiley and his co-workers Philippidis and Dewdney presented a full calculation on the explanation of the two-slit experiment in terms of Bohmian trajectories that arise for each particle moving under the influence of the quantum potential, resulting in the well-known interference patterns. Also the shift of the interference pattern which occurs in presence of a magnetic field in the Aharonov–Bohm effect could be explained as arising from the quantum potential. Relation to the measurement process The collapse of the wave function of the Copenhagen interpretation of quantum theory is explained in the quantum potential approach by the demonstration that, after a measurement, "all the packets of the multi-dimensional wave function that do not correspond to the actual result of measurement have no effect on the particle" from then on. Bohm and Hiley pointed out that Measurement then "involves a participatory transformation in which both the system under observation and the observing apparatus undergo a mutual participation so that the trajectories behave in a correlated manner, becoming correlated and separated into different, non-overlapping sets (which we call 'channels')". Quantum potential of an n-particle system The Schrödinger wave function of a many-particle quantum system cannot be represented in ordinary three-dimensional space. Rather, it is represented in configuration space, with three dimensions per particle. A single point in configuration space thus represents the configuration of the entire n-particle system as a whole. A two-particle wave function of identical particles of mass has the quantum potential where and refer to particle 1 and particle 2 respectively. This expression generalizes in straightforward manner to particles: In case the wave function of two or more particles is separable, then the system's total quantum potential becomes the sum of the quantum potentials of the two particles. Exact separability is extremely unphysical given that interactions between the system and its environment destroy the factorization; however, a wave function that is a superposition of several wave functions of approximately disjoint support will factorize approximately. Derivation for a separable quantum system That the wave function is separable means that factorizes in the form . Then it follows that also factorizes, and the system's total quantum potential becomes the sum of the quantum potentials of the two particles. In case the wave function is separable, that is, if factorizes in the form , the two one-particle systems behave independently. More generally, the quantum potential of an -particle system with separable wave function is the sum of quantum potentials, separating the system into independent one-particle systems. Formulation in terms of probability density Quantum potential in terms of the probability density function Bohm, as well as other physicists after him, have sought to provide evidence that the Born rule linking to the probability density function can be understood, in a pilot wave formulation, as not representing a basic law, but rather a theorem (called quantum equilibrium hypothesis) which applies when a quantum equilibrium is reached during the course of the time development under the Schrödinger equation. With Born's rule, and straightforward application of the chain and product rules the quantum potential, expressed in terms of the probability density function, becomes: Quantum force The quantum force , expressed in terms of the probability distribution, amounts to: Formulation in configuration space and in momentum space, as the result of projections M. R. Brown and B. Hiley showed that, as alternative to its formulation terms of configuration space (-space), the quantum potential can also be formulated in terms of momentum space (-space). In line with David Bohm's approach, Basil Hiley and mathematician Maurice de Gosson showed that the quantum potential can be seen as a consequence of a projection of an underlying structure, more specifically of a non-commutative algebraic structure, onto a subspace such as ordinary space (-space). In algebraic terms, the quantum potential can be seen as arising from the relation between implicate and explicate orders: if a non-commutative algebra is employed to describe the non-commutative structure of the quantum formalism, it turns out that it is impossible to define an underlying space, but that rather "shadow spaces" (homomorphic spaces) can be constructed and that in so doing the quantum potential appears. The quantum potential approach can be seen as a way to construct the shadow spaces. The quantum potential thus results as a distortion due to the projection of the underlying space into -space, in similar manner as a Mercator projection inevitably results in a distortion in a geographical map. There exists complete symmetry between the -representation, and the quantum potential as it appears in configuration space can be seen as arising from the dispersion of the momentum -representation. The approach has been applied to extended phase space, also in terms of a Duffin–Kemmer–Petiau algebra approach. Relation to other quantities and theories Relation to the Fisher information It can be shown that the mean value of the quantum potential is proportional to the probability density's Fisher information about the observable Using this definition for the Fisher information, we can write: Quantum potential as energy of internal motion associated with spin Giovanni Salesi, Erasmo Recami and co-workers showed in 1998 that, in agreement with the König's theorem, the quantum potential can be identified with the kinetic energy of the internal motion ("zitterbewegung") associated with the spin of a spin-1/2 particle observed in a center-of-mass frame. More specifically, they showed that the internal zitterbewegung velocity for a spinning, non-relativistic particle of constant spin with no precession, and in absence of an external field, has the squared value: from which the second term is shown to be of negligible size; then with it follows that Salesi gave further details on this work in 2009. In 1999, Salvatore Esposito generalized their result from spin-1/2 particles to particles of arbitrary spin, confirming the interpretation of the quantum potential as a kinetic energy for an internal motion. Esposito showed that (using the notation =1) the quantum potential can be written as: and that the causal interpretation of quantum mechanics can be reformulated in terms of a particle velocity where the "drift velocity" is and the "relative velocity" is , with and representing the spin direction of the particle. In this formulation, according to Esposito, quantum mechanics must necessarily be interpreted in probabilistic terms, for the reason that a system's initial motion condition cannot be exactly determined. Esposito explained that "the quantum effects present in the Schrödinger equation are due to the presence of a peculiar spatial direction associated with the particle that, assuming the isotropy of space, can be identified with the spin of the particle itself". Esposito generalized it from matter particles to gauge particles, in particular photons, for which he showed that, if modelled as , with probability function , they can be understood in a quantum potential approach. James R. Bogan, in 2002, published the derivation of a reciprocal transformation from the Hamilton-Jacobi equation of classical mechanics to the time-dependent Schrödinger equation of quantum mechanics which arises from a gauge transformation representing spin, under the simple requirement of conservation of probability. This spin-dependent transformation is a function of the quantum potential. Re-interpretation in terms of Clifford algebras B. Hiley and R. E. Callaghan re-interpret the role of the Bohm model and its notion of quantum potential in the framework of Clifford algebra, taking account of recent advances that include the work of David Hestenes on spacetime algebra. They show how, within a nested hierarchy of Clifford algebras , for each Clifford algebra an element of a minimal left ideal and an element of a right ideal representing its Clifford conjugation can be constructed, and from it the Clifford density element (CDE) , an element of the Clifford algebra which is isomorphic to the standard density matrix but independent of any specific representation. On this basis, bilinear invariants can be formed which represent properties of the system. Hiley and Callaghan distinguish bilinear invariants of a first kind, of which each stands for the expectation value of an element of the algebra which can be formed as , and bilinear invariants of a second kind which are constructed with derivatives and represent momentum and energy. Using these terms, they reconstruct the results of quantum mechanics without depending on a particular representation in terms of a wave function nor requiring reference to an external Hilbert space. Consistent with earlier results, the quantum potential of a non-relativistic particle with spin (Pauli particle) is shown to have an additional spin-dependent term, and the momentum of a relativistic particle with spin (Dirac particle) is shown to consist in a linear motion and a rotational part. The two dynamical equations governing the time evolution are re-interpreted as conservation equations. One of them stands for the conservation of energy; the other stands for the conservation of probability and of spin. The quantum potential plays the role of an internal energy which ensures the conservation of total energy. Relativistic and field-theoretic extensions Quantum potential and relativity Bohm and Hiley demonstrated that the non-locality of quantum theory can be understood as limit case of a purely local theory, provided the transmission of active information is allowed to be greater than the speed of light, and that this limit case yields approximations to both quantum theory and relativity. The quantum potential approach was extended by Hiley and co-workers to quantum field theory in Minkowski spacetime and to curved spacetime. Carlo Castro and Jorge Mahecha derived the Schrödinger equation from the Hamilton-Jacobi equation in conjunction with the continuity equation, and showed that the properties of the relativistic Bohm quantum potential in terms of the ensemble density can be described by the Weyl properties of space. In Riemann flat space, the Bohm potential is shown to equal the Weyl curvature. According to Castro and Mahecha, in the relativistic case, the quantum potential (using the d'Alembert operator and in the notation ) takes the form and the quantum force exerted by the relativistic quantum potential is shown to depend on the Weyl gauge potential and its derivatives. Furthermore, the relationship among Bohm's potential and the Weyl curvature in flat spacetime corresponds to a similar relationship among Fisher Information and Weyl geometry after introduction of a complex momentum. Diego L. Rapoport, on the other hand, associates the relativistic quantum potential with the metric scalar curvature (Riemann curvature). In relation to the Klein–Gordon equation for a particle with mass and charge, Peter R. Holland spoke in his book of 1993 of a "quantum potential-like term" that is proportional . He emphasized however that to give the Klein–Gordon theory a single-particle interpretation in terms of trajectories, as can be done for nonrelativistic Schrödinger quantum mechanics, would lead to unacceptable inconsistencies. For instance, wave functions that are solutions to the Klein–Gordon or the Dirac equation cannot be interpreted as the probability amplitude for a particle to be found in a given volume at time in accordance with the usual axioms of quantum mechanics, and similarly in the causal interpretation it cannot be interpreted as the probability for the particle to be in that volume at that time. Holland pointed out that, while efforts have been made to determine a Hermitian position operator that would allow an interpretation of configuration space quantum field theory, in particular using the Newton–Wigner localization approach, but that no connection with possibilities for an empirical determination of position in terms of a relativistic measurement theory or for a trajectory interpretation has so far been established. Yet according to Holland this does not mean that the trajectory concept is to be discarded from considerations of relativistic quantum mechanics. Hrvoje Nikolić derived as expression for the quantum potential, and he proposed a Lorentz-covariant formulation of the Bohmian interpretation of many-particle wave functions. He also developed a generalized relativistic-invariant probabilistic interpretation of quantum theory, in which is no longer a probability density in space but a probability density in space-time. Quantum potential in quantum field theory Starting from the space representation of the field coordinate, a causal interpretation of the Schrödinger picture of relativistic quantum theory has been constructed. The Schrödinger picture for a neutral, spin 0, massless field , with real-valued functionals, can be shown to lead to This has been called the superquantum potential by Bohm and his co-workers.<ref>Basil Hiley: The conceptual structure of the Bohm interpretation of quantum mechanics, Kalervo Vihtori Laurikainen et al (ed.): Symposium on the Foundations of Modern Physics 1994: 70 years of matter waves, Editions Frontières, , p. 99–117, p. 144</ref> Basil Hiley showed that the energy–momentum-relations in the Bohm model can be obtained directly from the energy–momentum tensor of quantum field theory and that the quantum potential is an energy term that is required for local energy–momentum conservation. He has also hinted that for particle with energies equal to or higher than the pair creation threshold, Bohm's model constitutes a many-particle theory that describes also pair creation and annihilation processes. Interpretation and naming of the quantum potential In his article of 1952, providing an alternative interpretation of quantum mechanics, Bohm already spoke of a "quantum-mechanical" potential. Bohm and Basil Hiley also called the quantum potential an information potential, given that it influences the form of processes and is itself shaped by the environment. Bohm indicated "The ship or aeroplane (with its automatic Pilot) is a self-active system, i.e. it has its own energy. But the form of its activity is determined by the information content concerning its environment that is carried by the radar waves. This is independent of the intensity of the waves. We can similarly regard the quantum potential as containing active information. It is potentially active everywhere, but actually active only where and when there is a particle." (italics in original). Hiley refers to the quantum potential as internal energy and as "a new quality of energy only playing a role in quantum processes". He explains that the quantum potential is a further energy term aside the well-known kinetic energy and the (classical) potential energy and that it is a nonlocal energy term that arises necessarily in view of the requirement of energy conservation; he added that much of the physics community's resistance against the notion of the quantum potential may have been due to scientists' expectations that energy should be local. Hiley has emphasized that the quantum potential, for Bohm, was "a key element in gaining insights into what could underlie the quantum formalism. Bohm was convinced by his deeper analysis of this aspect of the approach that the theory could not be mechanical. Rather, it is organic in the sense of Whitehead. Namely, that it was the whole that determined the properties of the individual particles and their relationship, not the other way round."	Physical sciences	Quantum mechanics	Physics
243074	https://en.wikipedia.org/wiki/Malic%20acid	Malic acid	Malic acid is an organic compound with the molecular formula . It is a dicarboxylic acid that is made by all living organisms, contributes to the sour taste of fruits, and is used as a food additive. Malic acid has two stereoisomeric forms (L- and D-enantiomers), though only the L-isomer exists naturally. The salts and esters of malic acid are known as malates. The malate anion is a metabolic intermediate in the citric acid cycle. Etymology The word 'malic' is derived from Latin , meaning 'apple'. The related Latin word , meaning 'apple tree', is used as the name of the genus Malus, which includes all apples and crabapples; and is the origin of other taxonomic classifications such as Maloideae, Malinae, and Maleae. Biochemistry L-Malic acid is the naturally occurring form, whereas a mixture of L- and D-malic acid is produced synthetically. Malate plays an important role in biochemistry. In the C4 carbon fixation process, malate is a source of CO2 in the Calvin cycle. In the citric acid cycle, (S)-malate is an intermediate, formed by the addition of an -OH group on the si face of fumarate. It can also be formed from pyruvate via anaplerotic reactions. Malate is also synthesized by the carboxylation of phosphoenolpyruvate in the guard cells of plant leaves. Malate, as a double anion, often accompanies potassium cations during the uptake of solutes into the guard cells in order to maintain electrical balance in the cell. The accumulation of these solutes within the guard cell decreases the solute potential, allowing water to enter the cell and promote aperture of the stomata. In food Malic acid was first isolated from apple juice by Carl Wilhelm Scheele in 1785. Antoine Lavoisier in 1787 proposed the name acide malique, which is derived from the Latin word for apple, mālum—as is its genus name Malus. In German it is named Äpfelsäure (or Apfelsäure) after plural or singular of a sour thing from the apple fruit, but the salt(s) are called Malat(e). Malic acid is the main acid in many fruits, including apricots, blackberries, blueberries, cherries, grapes, mirabelles, peaches, pears, plums, and quince, and is present in lower concentrations in other fruits, such as citrus. It contributes to the sourness of unripe apples. Sour apples contain high proportions of the acid. It is present in grapes and in most wines with concentrations sometimes as high as 5 g/L. It confers a tart taste to wine; the amount decreases with increasing fruit ripeness. The taste of malic acid is very clear and pure in rhubarb, a plant for which it is the primary flavor. It is also the compound responsible for the tart flavor of sumac spice. It is also a component of some artificial vinegar flavors, such as "salt and vinegar" flavored potato chips. The process of malolactic fermentation converts malic acid to much milder lactic acid. Malic acid occurs naturally in all fruits and many vegetables, and is generated in fruit metabolism. Malic acid, when added to food products, is denoted by E number E296. It is sometimes used with or in place of the less sour citric acid in sour sweets. These sweets are sometimes labeled with a warning stating that excessive consumption can cause irritation of the mouth. It is approved for use as a food additive in the EU, US and Australia and New Zealand (where it is listed by its INS number 296). Malic acid contains 10 kJ (2.39 kilocalories) of energy per gram. Production and main reactions Racemic malic acid is produced industrially by the double hydration of maleic anhydride. In 2000, American production capacity was 5,000 tons per year. The enantiomers may be separated by chiral resolution of the racemic mixture. S-Malic acid is obtained by fermentation of fumaric acid. Self-condensation of malic acid in the presence of fuming sulfuric acid gives the pyrone coumalic acid: Carbon monoxide and water are liberated during this reaction. Malic acid was important in the discovery of the Walden inversion and the Walden cycle, in which (−)-malic acid first is converted into (+)-chlorosuccinic acid by action of phosphorus pentachloride. Wet silver oxide then converts the chlorine compound to (+)-malic acid, which then reacts with PCl5 to the (−)-chlorosuccinic acid. The cycle is completed when silver oxide takes this compound back to (−)-malic acid. -malic acid is used to resolve α-phenylethylamine, a versatile resolving agent in its own right. Plant defense Soil supplementation with molasses increases microbial synthesis of malic acid. This is thought to occur naturally as part of soil microbe suppression of disease, so soil amendment with molasses can be used as a crop treatment in horticulture. Interactive pathway map	Physical sciences	Specific acids	Chemistry
243334	https://en.wikipedia.org/wiki/Green%27s%20theorem	Green's theorem	In vector calculus, Green's theorem relates a line integral around a simple closed curve to a double integral over the plane region (surface in ) bounded by . It is the two-dimensional special case of Stokes' theorem (surface in ). In one dimension, it is equivalent to the fundamental theorem of calculus. In three dimensions, it is equivalent to the divergence theorem. Theorem Let be a positively oriented, piecewise smooth, simple closed curve in a plane, and let be the region bounded by . If and are functions of defined on an open region containing and have continuous partial derivatives there, then where the path of integration along is counterclockwise. Application In physics, Green's theorem finds many applications. One is solving two-dimensional flow integrals, stating that the sum of fluid outflowing from a volume is equal to the total outflow summed about an enclosing area. In plane geometry, and in particular, area surveying, Green's theorem can be used to determine the area and centroid of plane figures solely by integrating over the perimeter. Proof when D is a simple region The following is a proof of half of the theorem for the simplified area D, a type I region where C1 and C3 are curves connected by vertical lines (possibly of zero length). A similar proof exists for the other half of the theorem when D is a type II region where C2 and C4 are curves connected by horizontal lines (again, possibly of zero length). Putting these two parts together, the theorem is thus proven for regions of type III (defined as regions which are both type I and type II). The general case can then be deduced from this special case by decomposing D into a set of type III regions. If it can be shown that and are true, then Green's theorem follows immediately for the region D. We can prove () easily for regions of type I, and () for regions of type II. Green's theorem then follows for regions of type III. Assume region D is a type I region and can thus be characterized, as pictured on the right, by where g1 and g2 are continuous functions on . Compute the double integral in (): Now compute the line integral in (). C can be rewritten as the union of four curves: C1, C2, C3, C4. With C1, use the parametric equations: x = x, y = g1(x), a ≤ x ≤ b. Then With C3, use the parametric equations: x = x, y = g2(x), a ≤ x ≤ b. Then The integral over C3 is negated because it goes in the negative direction from b to a, as C is oriented positively (anticlockwise). On C2 and C4, x remains constant, meaning Therefore, Combining () with (), we get () for regions of type I. A similar treatment yields () for regions of type II. Putting the two together, we get the result for regions of type III. Proof for rectifiable Jordan curves We are going to prove the following We need the following lemmas whose proofs can be found in: Now we are in position to prove the theorem: Proof of Theorem. Let be an arbitrary positive real number. By continuity of , and compactness of , given , there exists such that whenever two points of are less than apart, their images under are less than apart. For this , consider the decomposition given by the previous Lemma. We have Put . For each , the curve is a positively oriented square, for which Green's formula holds. Hence Every point of a border region is at a distance no greater than from . Thus, if is the union of all border regions, then ; hence , by Lemma 2. Notice that This yields We may as well choose so that the RHS of the last inequality is The remark in the beginning of this proof implies that the oscillations of and on every border region is at most . We have By Lemma 1(iii), Combining these, we finally get for some . Since this is true for every , we are done. Validity under different hypotheses The hypothesis of the last theorem are not the only ones under which Green's formula is true. Another common set of conditions is the following: The functions are still assumed to be continuous. However, we now require them to be Fréchet-differentiable at every point of . This implies the existence of all directional derivatives, in particular , where, as usual, is the canonical ordered basis of . In addition, we require the function to be Riemann-integrable over . As a corollary of this, we get the Cauchy Integral Theorem for rectifiable Jordan curves: Multiply-connected regions Theorem. Let be positively oriented rectifiable Jordan curves in satisfying where is the inner region of . Let Suppose and are continuous functions whose restriction to is Fréchet-differentiable. If the function is Riemann-integrable over , then Relationship to Stokes' theorem Green's theorem is a special case of the Kelvin–Stokes theorem, when applied to a region in the -plane. We can augment the two-dimensional field into a three-dimensional field with a z component that is always 0. Write F for the vector-valued function . Start with the left side of Green's theorem: The Kelvin–Stokes theorem: The surface is just the region in the plane , with the unit normal defined (by convention) to have a positive z component in order to match the "positive orientation" definitions for both theorems. The expression inside the integral becomes Thus we get the right side of Green's theorem Green's theorem is also a straightforward result of the general Stokes' theorem using differential forms and exterior derivatives: Relationship to the divergence theorem Considering only two-dimensional vector fields, Green's theorem is equivalent to the two-dimensional version of the divergence theorem: where is the divergence on the two-dimensional vector field , and is the outward-pointing unit normal vector on the boundary. To see this, consider the unit normal in the right side of the equation. Since in Green's theorem is a vector pointing tangential along the curve, and the curve C is the positively oriented (i.e. anticlockwise) curve along the boundary, an outward normal would be a vector which points 90° to the right of this; one choice would be . The length of this vector is So Start with the left side of Green's theorem: Applying the two-dimensional divergence theorem with , we get the right side of Green's theorem: Area calculation Green's theorem can be used to compute area by line integral. The area of a planar region is given by Choose and such that , the area is given by Possible formulas for the area of include History It is named after George Green, who stated a similar result in an 1828 paper titled An Essay on the Application of Mathematical Analysis to the Theories of Electricity and Magnetism. In 1846, Augustin-Louis Cauchy published a paper stating Green's theorem as the penultimate sentence. This is in fact the first printed version of Green's theorem in the form appearing in modern textbooks. George Green, An Essay on the Application of Mathematical Analysis to the Theories of Electricity and Magnetism (Nottingham, England: T. Wheelhouse, 1828). Green did not actually derive the form of "Green's theorem" which appears in this article; rather, he derived a form of the "divergence theorem", which appears on pages 10–12 of his Essay. In 1846, the form of "Green's theorem" which appears in this article was first published, without proof, in an article by Augustin Cauchy: A. Cauchy (1846) "Sur les intégrales qui s'étendent à tous les points d'une courbe fermée" (On integrals that extend over all of the points of a closed curve), Comptes rendus, 23: 251–255. (The equation appears at the bottom of page 254, where (S) denotes the line integral of a function k along the curve s that encloses the area S.) A proof of the theorem was finally provided in 1851 by Bernhard Riemann in his inaugural dissertation: Bernhard Riemann (1851) Grundlagen für eine allgemeine Theorie der Functionen einer veränderlichen complexen Grösse (Basis for a general theory of functions of a variable complex quantity), (Göttingen, (Germany): Adalbert Rente, 1867); see pages 8–9.	Mathematics	Multivariable and vector calculus	null
243447	https://en.wikipedia.org/wiki/Astronomical%20transit	Astronomical transit	In astronomy, a transit (or astronomical transit) is the passage of a celestial body directly between a larger body and the observer. As viewed from a particular vantage point, the transiting body appears to move across the face of the larger body, covering a small portion of it. The word "transit" refers to cases where the nearer object appears smaller than the more distant object. Cases where the nearer object appears larger and completely hides the more distant object are known as occultations. However, the probability of seeing a transiting planet is low because it is dependent on the alignment of the three objects in a nearly perfectly straight line. Many parameters of a planet and its parent star can be determined based on the transit. In the Solar System One type of transit involves the motion of a planet between a terrestrial observer and the Sun. This can happen only with inferior planets, namely Mercury and Venus (see transit of Mercury and transit of Venus). However, because a transit is dependent on the point of observation, the Earth itself transits the Sun if observed from Mars. In the solar transit by the Moon captured during calibration of the STEREO B spacecraft's ultraviolet imaging, the Moon appears much smaller than it does when seen from Earth, because the spacecraft–Moon separation was several times greater than the Earth–Moon distance. The term can also be used to describe the motion of a satellite across its parent planet, for instance one of the Galilean satellites (Io, Europa, Ganymede, Callisto) across Jupiter, as seen from Earth. Although rare, cases where four bodies are lined up do happen. One of these events occurred on 27 June 1586, when Mercury transited the Sun as seen from Venus at the same time as a transit of Mercury from Saturn and a transit of Venus from Saturn. Notable observations No missions were planned to coincide with the transit of Earth visible from Mars on 11 May 1984 and the Viking missions had been terminated a year previously. Consequently, the next opportunity to observe such an alignment will be in 2084. On 21 December 2012, the Cassini–Huygens probe, in orbit around Saturn, observed the planet Venus transiting the Sun. On 3 June 2014, the Mars rover Curiosity observed the planet Mercury transiting the Sun, marking the first time a planetary transit has been observed from a celestial body besides Earth. Mutual planetary transits In rare cases, one planet can pass in front of another. If the nearer planet appears smaller than the more distant one, the event is called a mutual planetary transit. Outside the Solar System The transit method can be used to discover exoplanets. As a planet eclipses/transits its host star it will block a portion of the light from the star. If the planet transits in-between the star and the observer the change in light can be measured to construct a light curve. Light curves are measured with a charge-coupled device. The light curve of a star can disclose several physical characteristics of the planet and star, such as density. Multiple transit events must be measured to determine the characteristics which tend to occur at regular intervals. Multiple planets orbiting the same host star can cause transit-timing variations (TTV). TTV is caused by the gravitational forces of all orbiting bodies acting upon each other. The probability of seeing a transit from Earth is low, however. The probability is given by the following equation. where Rstar and Rplanet are the radius of the star and planet, respectively, and a is the semi-major axis. Because of the low probability of a transit in any specific system, large selections of the sky must be regularly observed in order to see a transit. Hot Jupiters are more likely to be seen because of their larger radius and short semi-major axis. In order to find Earth-sized planets, red dwarf stars are observed because of their small radius. Even though transiting has a low probability it has proven itself to be a good technique for discovering exoplanets. In recent years, the discovery of extrasolar planets has prompted interest in the possibility of detecting their transits across their own stellar primaries. HD 209458b was the first such transiting planet to be detected. The transit of celestial objects is one of the few key phenomena used today for the study of exoplanetary systems. Today, transit photometry is the leading form of exoplanet discovery. As an exoplanet moves in front of its host star there is a dimming in the luminosity of the host star that can be measured. Larger planets make the dip in luminosity more noticeable and easier to detect. Followup observations using other methods are often carried out to ensure it is a planet. There are currently (December 2018) 2345 planets confirmed with Kepler light curves for stellar host. Contacts During a transit there are four "contacts", when the circumference of the small circle (small body disk) touches the circumference of the large circle (large body disk) at a single point. Historically, measuring the precise time of each point of contact was one of the most accurate ways to determine the positions of astronomical bodies. The contacts happen in the following order: First contact: the smaller body is entirely outside the larger body, moving inward ("exterior ingress") Second contact: the smaller body is entirely inside the larger body, moving further inward ("interior ingress") Third contact: the smaller body is entirely inside the larger body, moving outward ("interior egress") Fourth contact: the smaller body is entirely outside the larger body, moving outward ("exterior egress") A fifth named point is that of greatest transit, when the apparent centers of the two bodies are nearest to each other, halfway through the transit. Missions Since transit photometry allows for scanning large celestial areas with a simple procedure, it has been the most popular and successful form of finding exoplanets in the past decade and includes many projects, some of which have already been retired, others in use today, and some in progress of being planned and created. The most successful projects include HATNet, KELT, Kepler, and WASP, and some new and developmental stage missions such as TESS, HATPI, and others which can be found among the List of Exoplanet Search Projects. HATNet HATNet Project is a set of northern telescopes in Fred Lawrence Whipple Observatory, Arizona and Mauna Kea Observatories, HI, and southern telescopes around the globe, in Africa, Australia, and South America, under the HATSouth branch of the project. These are small aperture telescopes, just like KELT, and look at a wide field which allows them to scan a large area of the sky for possible transiting planets. In addition, their multitude and spread around the world allows for 24/7 observation of the sky so that more short-period transits can be caught. A third sub-project, HATPI, is currently under construction and will survey most of the night sky seen from its location in Chile. KELT KELT is a terrestrial telescope mission designed to search for transiting systems of planets of magnitude 8<M<10. It began operation in October 2004 in Winer Observatory and has a southern companion telescope added in 2009. KELT North observes "26-degree wide strip of sky that is overhead from North America during the year", while KELT South observes single target areas of the size 26 by 26 degrees. Both telescopes can detect and identify transit events as small as a 1% flux dip, which allows for detection of planetary systems similar to those in our planetary system. Kepler / K2 The Kepler space telescope served the Kepler mission between 7 March 2009 and 11 May 2013, where it observed one part of the sky in search of transiting planets within a 115 square degrees of the sky around the Cygnus, Lyra, and Draco constellations. After that, the satellite continued operating until 15 November 2018, this time changing its field along the ecliptic to a new area roughly every 75 days due to reaction wheel failure. TESS TESS was launched on 18 April 2018, and is planned to survey most of the sky by observing it strips defined along the right ascension lines for 27 days each. Each area surveyed is 27 by 90 degrees. Because of the positioning of sections, the area near TESS's rotational axis will be surveyed for up to 1 year, allowing for the identification of planetary systems with longer orbital periods.	Physical sciences	Celestial mechanics	Astronomy
243515	https://en.wikipedia.org/wiki/Community-supported%20agriculture	Community-supported agriculture	Community-supported agriculture (CSA model) or cropsharing is a system that connects producers and consumers within the food system closer by allowing the consumer to subscribe to the harvest of a certain farm or group of farms. It is an alternative socioeconomic model of agriculture and food distribution that allows the producer and consumer to share the risks of farming. The model is a subcategory of civic agriculture that has an overarching goal of strengthening a sense of community through local markets. Community-supported agriculture can be considered as a practice of Commoning. It is an example of community-led management of the production and distribution of goods and services. The organization of food provisioning through commoning is complementary to the horizontal axis of market mediated food provisioning and the verticality of the state distribution and regulation on food. As a model where market agents do not interact solely as competitors but as “members of a community collaborating in pursuing a collective action for the commonwealth” it is also recognized and supported by public policies in some countries. Such frameworks of collaboration between public administration and the cooperative sector are known as Public-Commons-Partnerships (PCP) and have also been established in relation to food. As a prefigurative practice that decommodifies food and “strengthens the imaginary of community as a source of reward and space of emancipation“ CSA has been acknowledged as an important step-stone in a sustainability transition in agri-food systems. In return for subscribing to a harvest, subscribers receive either a weekly or bi-weekly box of produce or other farm goods. This includes in-season fruits, vegetables, and can expand to dried goods, eggs, milk, meat, etc. Typically, farmers try to cultivate a relationship with subscribers by sending weekly letters of what is happening on the farm, inviting them for harvest, or holding an open-farm event. Some CSAs provide for contributions of labor in lieu of a portion of subscription costs. The term CSA is mostly used in the United States, Canada and the UK but a variety of similar production and economic sub-systems are in use worldwide and in Austria and Germany as Solidarische Landwirtschaft (, abbreviated to Solawi). History The term "community-supported agriculture" was coined in the northeastern United States in the 1980s, influenced by European biodynamic agriculture ideas formulated by Rudolf Steiner. Two European farmers, Jan Vander Tuin from Switzerland and Trauger Groh from Germany, brought European biodynamic farming ideas to the United States in the mid-1980s. Vander Tuin had co-founded a community-supported agricultural project named Topinambur located near Zurich, Switzerland. Coinage of the term "community-supported agriculture" stems from Vander Tuin. This influence led to the separate and simultaneous creation of two CSAs in 1986. The CSA Garden at Great Barrington was created in Massachusetts by Jan Vander Tuin, Susan Witt, and Robyn Van En. The Temple-Wilton Community Farm was created in New Hampshire by Anthony Graham, Trauger Groh, and Lincoln Geiger. The CSA Garden at Great Barrington remained together until 1990 when many members left to form the Mahaiwe Harvest CSA. One of the original founders, Robyn Van En, became incredibly influential in the CSA movement in America and founded CSA North America in 1992. The Temple-Wilton Community Garden was more successful and still operates as a CSA today. It became an important member of the Wilton community and it receives funding from state, federal, and local sources. A parallel model called Teikei existed in Japan as early as the mid-1960s. Similarly, Dr. Booker T. Whatley, a professor of agriculture in Alabama, advocated for Clientele Membership Clubs as early as the 1960s. Since the 1980s, community supported farms have been organized throughout North America—mainly in New England, the Northwest, the Pacific coast, the Upper-Midwest and Canada. North America now has at least 13,000 CSA farms of which 12,549 are in the US according to the United States Department of Agriculture in 2007. The rise of CSAs seems to be correlated with the increase in awareness of the environmental movement in the United States, and because food acquisition from local sources can reduce greenhouse gas emissions, CSAs contribute to climate change mitigation. CSAs have even become popular in urban environments such as the New York City Coalition Against Hunger's CSA program that helps serve under-served communities. One of the largest subscription CSAs was Capay Inc. in Capay Valley, California which in 2010 delivered boxes to 13,000 customers a week as well as selling at 15 farmers markets, operating a retail store, and delivering special orders to restaurants. Urgenci, based in France, helps network together consumers and producers across Europe, the Mediterranean, and West Africa. CSA was introduced to China following a series of food safety scandals in the late 2000s. It was estimated that there were more than 500 CSA farms in China by 2017. They have been a critical force in the development of the organic and ecological farming in China. Chinese CSA farmers, researchers and civil society organizations gather annually at the national CSA symposium held since 2009. Much of the growth in women labour participation in agriculture is outside the "male dominated field of conventional agriculture". In community supported agriculture women represent 40 percent of farm operators. International Even if the systems of community-supported agriculture vary in different countries, there are a number of umbrella-organizations connecting the farms. In the United States the governmental program SARE offered grants for research and education projects that advance sustainable agricultural practices like CSA. In Germany and Austria the CSA groups founded the Bundesnetzwerk Solidarische Landwirtschaft (Federal Network of CSA-farms) in 2011. Switzerland In Switzerland, community farming is often referred to as Solidarische Landwirtschaft (, abbreviated to Solawi). Among others there are two important forms of organization for CSA in Switzerland: the cooperative and the individual initiative. Food cooperatives are the oldest form of organization and the first CSA projects in Switzerland were organized as cooperatives. They belong to the category of farmer-shareholder cooperatives. Cooperatives are based on direct cooperation between farmers and consumers. Consumers or cooperative members actively participate in the management of agricultural production together with the producers. Producers assume the day-to-day management and work on the farm and are often formally employed by the cooperative, as well as being part of it. Cooperative members participate in production and distribution costs by purchasing shares in the cooperative and paying annual fees for the delivery of vegetables. Cooperative members also participate in some production or distribution tasks as part of their commitment to the cooperative. Some CSAs are also initiated by producers, often created to open a new distribution channel. This form of organization is often referred to as farmer-run CSAs. Producers offer baskets with the farm's produce to any interested consumer (consumers are not organized, but have individual contracts with the farmer).The baskets with vegetables and/or fruits are delivered to consumers on a regular basis. The distribution and management of production is the responsibility of the producer and there is usually no collaboration or shared investment between consumers and producers. For the producer it means a capacity and possibility to diversify production and has the advantage of a new distribution channel with very low entry costs. Regional networks in Switzerland The Fédération Romande d'Agriculture Contractuelle de Proximité (FRACP) is a French-speaking federation of CSAs and supports the creation of new CSAs. It was created in 2008 with 13 CSA groups as founding members. FRACP is sponsored by Uniterre, a small farmers' union that is part of La Via Campesina and promotes the concept of food sovereignty. Verband Regionale Vertragslandwirtschaft (RVL) is a German-speaking organization created in 2011 with the help of FRACP. RVL collaborates with the Kooperationsstelle für solidarische Landwirtschaft and offers courses on community supported agriculture. Italy The CSA (Community-Supported Agriculture) model in Italy is a relatively young movement that began to gain traction in 2011. The first CSA established in Italy was the C.A.P.S. (Agrarian Community of Social Promotion) in Pisa, while the largest CSA is Arvaia in Bologna, boasting 220 active members and 500 associates. As of 2021, a survey conducted by Numes, a project born in collaboration with the Arvaia CSA, identified 15 formal CSAs, although the actual number is likely higher. It is closely connected to the history of GAS (Gruppi di acquisto solidale) – Solidarity Purchase Groups). GAS and CSA share similar ethical values and organizational structures, operating based on principles of solidarity, mutuality, and sustainability. However, there are distinct differences between CSAs and GASs in terms of risk-sharing. In a CSA, members choose to provide financial support to farmers, thereby sharing the risks inherent in agricultural work. On the other hand, GAS members do not enter into a formal contract that obligates them to share any potential costs. United Kingdom The UK's first CSA was established in Findhorn, Scotland in 1994. The national umbrella organisation was set up in 2013 - this is the Community Supported Agriculture Network . This network aims to help people set up their CSAs, as well as providing helpful information to help these relations to run more successfully. Registered CSAs are mostly distributed across the east of Scotland. Yet, it is likely there are many more ‘informal’ CSAs on the west coast and the Isles. Across Wales and Northern Ireland, CSA's are spread evenly. In England, most CSA's are located in the South East. CSA enterprises across the UK have been growing fast, with most of the farm to community connections beginning and remaining as grassroots initiates, despite limited funding from government or private sector. Funding Support came from the Soil Association’s programme within the Big Lottery funded “Making Local Food Work” scheme which ran from 2007 to 2012. This supported the growth of CSAs across each of the countries, through providing funding that enabled the use of community and social enterprise approaches to link consumers and producers in the local food-related third sector. Socio-economic model CSAs create direct connections between producers and consumers through alternative markets and the members and farmers share the risk of farming. The goals of the first CSA model in the US were to have the producer and consumer to come into the market as equals and make an exchange with fair prices and fair wages. The consumer pays for things such as transparency, environmental stewardship, producer relationships, etc. The farmers engaged in CSAs do so to fulfill goals other than income and are not compensated fairly in these exchanges. This kind of market holds "economic rents" where the consumer surplus comes from the consumers' willingness to pay for something further than the product as well as for the products inputs themselves. Although these markets still exist within a larger capitalist economy, they are able to exist because of the "economic rents" that are collected. CSA system CSAs generally focus on the production of high quality foods for a local community, often using organic or biodynamic farming methods, and a shared risk membership–marketing structure. This kind of farming operates with a much greater degree of involvement of consumers and other stakeholders than usual—resulting in a stronger consumer-producer relationship. The core design includes developing a cohesive consumer group that is willing to fund a whole season's budget in order to get quality foods. The system has many variations on how the farm budget is supported by the consumers and how the producers then deliver the foods. CSA theory purports that the more a farm embraces whole-farm, whole-budget support, the more it can focus on quality and reduce the risk of food waste. Structure Community-supported agriculture farms in the United States today share three common characteristics: an emphasis on community and/or local produce, share or subscriptions sold prior to season, and weekly deliveries to members/subscribers. Though CSA operation varies from farm to farm and has evolved over time, these three characteristics have remained constant. The functioning of a CSA also relies on four practical arrangements: for farmers to know the needs of a community, for consumers to have the opportunity to express to farmers what their needs and financial limitations are, for commitments between farmers and consumers to be consciously established, and for farmers needs to be recognized. From this base, four main types of CSAs have been developed: Farmer managed: A farmer sets up and maintains a CSA, recruits subscribers, and controls management of the CSA. Shareholder/subscriber: Local residents set up a CSA and hire a farmer to grow crops, and shareholders/subscribers control most management. Farmer cooperative: Multiple farmers develop a CSA program. Farmer-shareholder cooperative: Farmers and local residents set up and cooperatively manage a CSA. In most original CSAs, a core group of members existed. This core group of members helped to make decisions about and run the CSA including marketing, distribution, administrative, and community organization functions. CSAs with a core group of members are most profitable and successful. However, in 1999, 72 percent of CSAs did not have a core group of members. CSAs with a core group of members operate more successfully as a farmer-shareholder cooperative and CSAs without a group of core members rely much more on subscriptions and run most prominently as shareholder/subscriber CSAs. Ideology Community-supported agriculture in America was influenced by the ideas of Rudolf Steiner, an Austrian philosopher. He developed the concepts of anthroposophy and biodynamic agriculture. The Temple-Wilton Community Farm used his ideas to develop three main goals of CSAs: New forms of property ownership: the idea that land should be held in common by a community through a legal trust, which leases the land to farmers New forms of cooperation: the idea that a network of human relationships should replace the traditional system of employers and employees New forms of economy: that the economy should not be based on increasing profit, but should be based on the actual needs of the people and land involved in an enterprise As CSAs have increased in both number and size since they were first developed, they have also changed ideologically. While original CSAs and some more current CSAs are still philosophically oriented, most CSAs today are commercially oriented and community-supported agriculture is predominantly seen as a beneficial marketing strategy. This has led to three ideologically based types of CSAs. The first type is instrumental, the CSA is considered a market in the traditional sense, instead of an alternative form of economy and relationship. The second type is functional; there is a relationship of solidarity between the farmer and the subscribers, but this extends mostly to social functions, not managerial or administrative functions. This is the most common type of CSA. The final type is collaborative; this is the closest to the original aims of CSAs where the relationship between the farmer and the subscribers is seen as a partnership. Distribution and marketing methods Shares of a CSA originally and predominantly consist of produce. In more recent years, shares have diversified and include non-produce products including eggs, meat, flowers, honey, dairy and soaps. Share prices vary from CSA to CSA. Shares are sold as full shares, which feed 2 to 5 people, and half shares, which feed 1 to 3 people. Prices range from $200 to $500 per season. Full shares are sold at a median of $400 and half shares are sold at a median of $250. Share prices are mostly determined by overhead costs of production, but are also determined by share prices of other CSAs, variable costs of production, market forces, and income level of the community. Many CSAs have payment plans and low-income options. Shares are distributed in several different ways. Shares are most often distributed weekly. Most CSAs allow share pick up at the farm. Shares are also distributed through regional dropoff, direct home or office dropoff, farmers' markets, and community center/church dropoff. For example, the new "Farmie Markets" of upstate New York take orders online and have a number of farmers who send that week's orders to a central point in a limited region, for distribution by the organizers. CSAs market their farms and shares in different ways. CSAs employ different channels of marketing to diversify their sales efforts and increase subscriptions. CSAs use local farmers' markets, restaurants, on-farm retail, wholesale to natural food stores, and wholesale to local groceries in addition to their CSAs to market shares. One problem that CSAs encounter is over-production, so CSAs often sell their produce and products in ways other than shares. Often, CSA farms also sell their products at local farmers' markets. Excess products are sometimes given to food banks. Challenges for farmers Many CSA farmers can capitalize on a closer relationship between customers and their food, since some customers will pay more (an economic rent if this puts the price above the cost of production) if they know where it is coming from, who is involved, and have special access to it. However, some farmers participating in community-supported agriculture do not experience the economic benefits that they are perceived to obtain by participating in an alternative community-based arrangement. Galt's 2013 study of CSA farmers found that many farmers charged lower fees and prices for their goods than would provide them with financial security. This study suggested that farmers may charge less than they need to earn fair wages due to undervaluing their expenses and to offset the high costs of CSA products and make it more affordable for customers; see moral economy.	Technology	Agriculture, labor and economy	null
243577	https://en.wikipedia.org/wiki/Jaw	Jaw	The jaws are a pair of opposable articulated structures at the entrance of the mouth, typically used for grasping and manipulating food. The term jaws is also broadly applied to the whole of the structures constituting the vault of the mouth and serving to open and close it and is part of the body plan of humans and most animals. Arthropods In arthropods, the jaws are chitinous and oppose laterally, and may consist of mandibles or chelicerae. These jaws are often composed of numerous mouthparts. Their function is fundamentally for food acquisition, conveyance to the mouth, and/or initial processing (mastication or chewing). Many mouthparts and associate structures (such as pedipalps) are modified legs. Vertebrates In most vertebrates, the jaws are bony or cartilaginous and oppose vertically, comprising an upper jaw and a lower jaw. The vertebrate jaw is derived from the most anterior two pharyngeal arches supporting the gills, and usually bears numerous teeth. Fish The vertebrate jaw probably originally evolved in the Silurian period and appeared in the Placoderm fish which further diversified in the Devonian. The two most anterior pharyngeal arches are thought to have become the jaw itself and the hyoid arch, respectively. The hyoid system suspends the jaw from the braincase of the skull, permitting great mobility of the jaws. While there is no fossil evidence directly to support this theory, it makes sense in light of the numbers of pharyngeal arches that are visible in extant jawed vertebrates (the Gnathostomes), which have seven arches, and primitive jawless vertebrates (the Agnatha), which have nine. The original selective advantage offered by the jaw may not be related to feeding, but rather to increased respiration efficiency. The jaws were used in the buccal pump (observable in modern fish and amphibians) that pumps water across the gills of fish or air into the lungs in the case of amphibians. Over evolutionary time the more familiar use of jaws (to humans), in feeding, was selected for and became a very important function in vertebrates. Many teleost fish have substantially modified jaws for suction feeding and jaw protrusion, resulting in highly complex jaws with dozens of bones involved. Amphibians, reptiles, and birds The jaw in tetrapods is substantially simplified compared to fish. Most of the upper jaw bones (premaxilla, maxilla, jugal, quadratojugal, and quadrate) have been fused to the braincase, while the lower jaw bones (dentary, splenial, angular, surangular, and articular) have been fused together into a unit called the mandible. The jaw articulates via a hinge joint between the quadrate and articular. The jaws of tetrapods exhibit varying degrees of mobility between jaw bones. Some species have jaw bones completely fused, while others may have joints allowing for mobility of the dentary, quadrate, or maxilla. The snake skull shows the greatest degree of cranial kinesis, which allows the snake to swallow large prey items. Mammals In mammals, the jaws are made up of the mandible (lower jaw) and the maxilla (upper jaw). In the ape, there is a reinforcement to the lower jaw bone called the simian shelf. In the evolution of the mammalian jaw, two of the bones of the jaw structure (the articular bone of the lower jaw, and quadrate) were reduced in size and incorporated into the ear, while many others have been fused together. As a result, mammals show little or no cranial kinesis, and the mandible is attached to the temporal bone by the temporomandibular joints. Temporomandibular joint dysfunction is a common disorder of these joints, characterized by pain, clicking and limitation of mandibular movement. Especially in the therian mammal, the premaxilla that constituted the anterior tip of the upper jaw in reptiles has reduced in size; and most of the mesenchyme at the ancestral upper jaw tip has become a protruded mammalian nose. Sea urchins Sea urchins possess unique jaws which display five-part symmetry, termed the Aristotle's lantern. Each unit of the jaw holds a single, perpetually growing tooth composed of crystalline calcium carbonate.	Biology and health sciences	Skeletal system	null
243613	https://en.wikipedia.org/wiki/Network%20interface%20controller	Network interface controller	A network interface controller (NIC, also known as a network interface card, network adapter, LAN adapter and physical network interface) is a computer hardware component that connects a computer to a computer network. Early network interface controllers were commonly implemented on expansion cards that plugged into a computer bus. The low cost and ubiquity of the Ethernet standard means that most newer computers have a network interface built into the motherboard, or is contained into a USB-connected dongle. Modern network interface controllers offer advanced features such as interrupt and DMA interfaces to the host processors, support for multiple receive and transmit queues, partitioning into multiple logical interfaces, and on-controller network traffic processing such as the TCP offload engine. Purpose The network controller implements the electronic circuitry required to communicate using a specific physical layer and data link layer standard such as Ethernet or Wi-Fi. This provides a base for a full network protocol stack, allowing communication among computers on the same local area network (LAN) and large-scale network communications through routable protocols, such as Internet Protocol (IP). The NIC allows computers to communicate over a computer network, either by using cables or wirelessly. The NIC is both a physical layer and data link layer device, as it provides physical access to a networking medium and, for IEEE 802 and similar networks, provides a low-level addressing system through the use of MAC addresses that are uniquely assigned to network interfaces. Implementation Network controllers were originally implemented as expansion cards that plugged into a computer bus. The low cost and ubiquity of the Ethernet standard means that most new computers have a network interface controller built into the motherboard. Newer server motherboards may have multiple network interfaces built-in. The Ethernet capabilities are either integrated into the motherboard chipset or implemented via a low-cost dedicated Ethernet chip. A separate network card is typically no longer required unless additional independent network connections are needed or some non-Ethernet type of network is used. A general trend in computer hardware is towards integrating the various components of systems on a chip, and this is also applied to network interface cards. An Ethernet network controller typically has an 8P8C socket where the network cable is connected. Older NICs also supplied BNC, or AUI connections. Ethernet network controllers typically support 10 Mbit/s Ethernet, 100 Mbit/s Ethernet, and Ethernet varieties. Such controllers are designated as 10/100/1000, meaning that they can support data rates of 10, 100 or . 10 Gigabit Ethernet NICs are also available, and, , are beginning to be available on computer motherboards. Modular designs like SFP and SFP+ are highly popular, especially for fiber-optic communication. These define a standard receptacle for media-dependent transceivers, so users can easily adapt the network interface to their needs. LEDs adjacent to or integrated into the network connector inform the user of whether the network is connected, and when data activity occurs. The NIC may include ROM to store its factory-assigned MAC address. The NIC may use one or more of the following techniques to indicate the availability of packets to transfer: Polling is where the CPU examines the status of the peripheral under program control. Interrupt-driven I/O is where the peripheral alerts the CPU that it is ready to transfer data. NICs may use one or more of the following techniques to transfer packet data: Programmed input/output, where the CPU moves the data to or from the NIC to memory. Direct memory access (DMA), where a device other than the CPU assumes control of the system bus to move data to or from the NIC to memory. This removes load from the CPU but requires more logic on the card. In addition, a packet buffer on the NIC may not be required and latency can be reduced. Performance and advanced functionality Multiqueue NICs provide multiple transmit and receive queues, allowing packets received by the NIC to be assigned to one of its receive queues. The NIC may distribute incoming traffic between the receive queues using a hash function. Each receive queue is assigned to a separate interrupt; by routing each of those interrupts to different CPUs or CPU cores, processing of the interrupt requests triggered by the network traffic received by a single NIC can be distributed improving performance. The hardware-based distribution of the interrupts, described above, is referred to as receive-side scaling (RSS). Purely software implementations also exist, such as the receive packet steering (RPS), receive flow steering (RFS), and Intel Flow Director. Further performance improvements can be achieved by routing the interrupt requests to the CPUs or cores executing the applications that are the ultimate destinations for network packets that generated the interrupts. This technique improves locality of reference and results in higher overall performance, reduced latency and better hardware utilization because of the higher utilization of CPU caches and fewer required context switches. With multi-queue NICs, additional performance improvements can be achieved by distributing outgoing traffic among different transmit queues. By assigning different transmit queues to different CPUs or CPU cores, internal operating system contentions can be avoided. This approach is usually referred to as transmit packet steering (XPS). Some products feature NIC partitioning (NPAR, also known as port partitioning) that uses SR-IOV virtualization to divide a single 10 Gigabit Ethernet NIC into multiple discrete virtual NICs with dedicated bandwidth, which are presented to the firmware and operating system as separate PCI device functions. Some NICs provide a TCP offload engine to offload processing of the entire TCP/IP stack to the network controller. It is primarily used with high-speed network interfaces, such as Gigabit Ethernet and 10 Gigabit Ethernet, for which the processing overhead of the network stack becomes significant. Some NICs offer integrated field-programmable gate arrays (FPGAs) for user-programmable processing of network traffic before it reaches the host computer, allowing for significantly reduced latencies in time-sensitive workloads. Moreover, some NICs offer complete low-latency TCP/IP stacks running on integrated FPGAs in combination with userspace libraries that intercept networking operations usually performed by the operating system kernel; Solarflare's open-source OpenOnload network stack that runs on Linux is an example. This kind of functionality is usually referred to as user-level networking.	Technology	Networks	null
243709	https://en.wikipedia.org/wiki/Cubic%20function	Cubic function	In mathematics, a cubic function is a function of the form that is, a polynomial function of degree three. In many texts, the coefficients , , , and are supposed to be real numbers, and the function is considered as a real function that maps real numbers to real numbers or as a complex function that maps complex numbers to complex numbers. In other cases, the coefficients may be complex numbers, and the function is a complex function that has the set of the complex numbers as its codomain, even when the domain is restricted to the real numbers. Setting produces a cubic equation of the form whose solutions are called roots of the function. The derivative of a cubic function is a quadratic function. A cubic function with real coefficients has either one or three real roots (which may not be distinct); all odd-degree polynomials with real coefficients have at least one real root. The graph of a cubic function always has a single inflection point. It may have two critical points, a local minimum and a local maximum. Otherwise, a cubic function is monotonic. The graph of a cubic function is symmetric with respect to its inflection point; that is, it is invariant under a rotation of a half turn around this point. Up to an affine transformation, there are only three possible graphs for cubic functions. Cubic functions are fundamental for cubic interpolation. History Critical and inflection points The critical points of a cubic function are its stationary points, that is the points where the slope of the function is zero. Thus the critical points of a cubic function defined by , occur at values of such that the derivative of the cubic function is zero. The solutions of this equation are the -values of the critical points and are given, using the quadratic formula, by The sign of the expression inside the square root determines the number of critical points. If it is positive, then there are two critical points, one is a local maximum, and the other is a local minimum. If , then there is only one critical point, which is an inflection point. If , then there are no (real) critical points. In the two latter cases, that is, if is nonpositive, the cubic function is strictly monotonic. See the figure for an example of the case . The inflection point of a function is where that function changes concavity. An inflection point occurs when the second derivative is zero, and the third derivative is nonzero. Thus a cubic function has always a single inflection point, which occurs at Classification The graph of a cubic function is a cubic curve, though many cubic curves are not graphs of functions. Although cubic functions depend on four parameters, their graph can have only very few shapes. In fact, the graph of a cubic function is always similar to the graph of a function of the form This similarity can be built as the composition of translations parallel to the coordinates axes, a homothecy (uniform scaling), and, possibly, a reflection (mirror image) with respect to the -axis. A further non-uniform scaling can transform the graph into the graph of one among the three cubic functions This means that there are only three graphs of cubic functions up to an affine transformation. The above geometric transformations can be built in the following way, when starting from a general cubic function Firstly, if , the change of variable allows supposing . After this change of variable, the new graph is the mirror image of the previous one, with respect of the -axis. Then, the change of variable provides a function of the form This corresponds to a translation parallel to the -axis. The change of variable corresponds to a translation with respect to the -axis, and gives a function of the form The change of variable corresponds to a uniform scaling, and give, after multiplication by a function of the form which is the simplest form that can be obtained by a similarity. Then, if , the non-uniform scaling gives, after division by where has the value 1 or –1, depending on the sign of . If one defines the latter form of the function applies to all cases (with and ). Symmetry For a cubic function of the form the inflection point is thus the origin. As such a function is an odd function, its graph is symmetric with respect to the inflection point, and invariant under a rotation of a half turn around the inflection point. As these properties are invariant by similarity, the following is true for all cubic functions. The graph of a cubic function is symmetric with respect to its inflection point, and is invariant under a rotation of a half turn around the inflection point. Collinearities The tangent lines to the graph of a cubic function at three collinear points intercept the cubic again at collinear points. This can be seen as follows. As this property is invariant under a rigid motion, one may suppose that the function has the form If is a real number, then the tangent to the graph of at the point is the line . So, the intersection point between this line and the graph of can be obtained solving the equation , that is which can be rewritten and factorized as So, the tangent intercepts the cubic at So, the function that maps a point of the graph to the other point where the tangent intercepts the graph is This is an affine transformation that transforms collinear points into collinear points. This proves the claimed result. Cubic interpolation Given the values of a function and its derivative at two points, there is exactly one cubic function that has the same four values, which is called a cubic Hermite spline. There are two standard ways for using this fact. Firstly, if one knows, for example by physical measurement, the values of a function and its derivative at some sampling points, one can interpolate the function with a continuously differentiable function, which is a piecewise cubic function. If the value of a function is known at several points, cubic interpolation consists in approximating the function by a continuously differentiable function, which is piecewise cubic. For having a uniquely defined interpolation, two more constraints must be added, such as the values of the derivatives at the endpoints, or a zero curvature at the endpoints.	Mathematics	Specific functions	null
243718	https://en.wikipedia.org/wiki/Flashlight	Flashlight	A flashlight (US English) or electric torch (Commonwealth English), usually shortened to torch, is a portable hand-held electric lamp. Formerly, the light source typically was a miniature incandescent light bulb, but these have been displaced by light-emitting diodes (LEDs) since the early 2000s. A typical flashlight consists of the light source mounted in a reflector, a transparent cover (sometimes combined with a lens) to protect the light source and reflector, a battery, and a switch, all enclosed in a case. The invention of the dry cell and miniature incandescent electric lamps made the first battery-powered flashlights possible around 1899. Today, flashlights use mostly light-emitting diodes and run on disposable or rechargeable batteries. Some are powered by the user turning a crank, shaking the lamp, or squeezing it. Some have solar panels to recharge the battery. Flashlights are used as a light source outdoors, in places without permanently installed lighting, during power outages, or when a portable light source is needed. In addition to the general-purpose, hand-held flashlight, many forms have been adapted for special uses. Head- or helmet-mounted flashlights designed for miners and campers leave both hands free. Some flashlights can be used under water or in flammable atmospheres. Etymology Early flashlights ran on zinc–carbon batteries, which could not provide a steady electric current and required periodic "rest" to continue functioning. Because these early flashlights also used energy-inefficient carbon-filament bulbs, "resting" occurred at short intervals. Consequently, they could be used only in brief flashes, hence the common North American name "flashlight". History The first dry cell battery was invented in 1887. Unlike previous batteries, it used a paste electrolyte instead of a liquid. This was the first battery suitable for portable electrical devices, as it did not spill or break easily and worked in any orientation. The first mass-produced dry cell batteries came in 1896, and the invention of portable electric lights soon followed. Portable hand-held electric lights offered advantages in convenience and safety over (combustion) torches, candles and lanterns. The electric lamp was odorless, smokeless, and emitted less heat than combustion-powered lighting. It could be instantly turned on and off, and avoided fire risk. On January 10, 1899, British inventor David Misell obtained U.S. Patent No. 617,592, assigned to American Electrical Novelty and Manufacturing Company. This "electric device" designed by Misell was powered by "D" batteries laid front to back in a paper tube with the light bulb and a rough brass reflector at the end. The company donated some of these devices to the New York City police, who responded favorably to them. Carbon-filament bulbs and fairly crude dry cells made early flashlights an expensive novelty, with low sales and low manufacturer interest. Development of the tungsten-filament lamp in 1904, with three times the efficacy of carbon filament types, along with improved batteries in varying sizes made flashlights more useful and popular. The advantage of instant control, and the absence of flame, meant that hand-held electric lights began to replace combustion-based lamps such as the hurricane lantern. By 1907, several types of flashlights were available: the tubular hand-held variety, a lantern style that could be set down for extended use, pocket-size penlights for close work, and large reflector searchlight-type lamps for lighting distant objects. In 1922 there were an estimated 10 million flashlight users in the United States, with annual sales of renewal batteries and flashlights at $20 million, comparable to sales of many line-operated electrical appliances. Flashlights became very popular in China; by the end of the 1930s, 60 companies made flashlights, some selling for as little as one-third the cost of equivalent imported models. Miniature lamps developed for flashlight and automotive uses became an important sector of the incandescent lamp manufacturing business. LED flashlights were introduced in the early 2000s. Maglite made their first LED flashlight in 2006. Incandescent Incandescent flashlights use incandescent light bulbs, which consists of a glass bulb and a tungsten filament. The bulbs are under vacuum or filled with argon, krypton, or xenon. Some high-power incandescent flashlights use a halogen lamp where the bulb contains a halogen gas such as iodine or bromine to improve the life and efficacy of the bulb. In all but disposable or novelty flashlights, the bulb is user-replaceable; the bulb life may be only a few hours. The light output of an incandescent lamp in a flashlight varies widely depending on the type of lamp. A miniature keychain lamp produces one or two lumens. A two-D-cell flashlight using a common prefocus-style miniature lamp produces on the order of 15 to 20 lumens of light and a beam of about 200 candlepower. One popular make of rechargeable focusing flashlight uses a halogen lamp and produces 218 lumens. By comparison, a 60-watt household incandescent lamp will produce about 900 lumens. The luminous efficacy or lumens produced per watt of input of flashlight bulbs varies over the approximate range of 8 to 22 lumens/watt, depending on the size of the bulb and the fill gas, with halogen-filled 12-volt lamps having the highest efficiency. LED Powerful white-light-emitting diodes (LEDs) have mostly replaced incandescent bulbs in practical flashlights. LEDs existed for decades, mainly as low-power indicator lights. In 1999, Lumileds Corporation of San Jose, California, introduced the Luxeon LED, a high-power white-light emitter. This made possible LED flashlights with lower power consumption and running time better than incandescent flashlights with similar light output. The first Luxeon LED flashlight was the Arc LS, designed in 2001. White LEDs in 5 mm diameter packages produce only a few lumens each; many units may be grouped together to provide additional light. Higher-power LEDs, drawing more than 100 milliamperes each, simplify the optical design problem of producing a powerful and tightly controlled beam. LEDs can be significantly more efficient than incandescent lamps, with white LEDs producing on the order of 100 lumens for every watt, compared to 8-10 lumens per watt of small incandescent bulbs. An LED flashlight has a longer battery life than an incandescent flashlight with comparable output. LEDs are also less fragile than glass lamps. LED lamps have different spectra of light compared to incandescent sources, and are made in several ranges of color temperature and color rendering index. Since the LED has a long life compared to the usual life of a flashlight, very often it is permanently installed. Flashlights made for an incandescent lamp can often be upgraded to a more efficient LED lamp. LEDs generally must have some kind of control to limit current through the diode. Flashlights using one or two disposable 1.5-volt cells require a boost converter to provide the higher voltage required by a white LED, which needs around 3.4 volts to function. Flashlights using three or more dry cells may only use a resistor to limit current. Some flashlights electronically regulate the current through the LEDs to stabilize light output as the batteries discharge. LEDs maintain nearly constant color temperature regardless of input voltage or current, while the color temperature of an incandescent bulb rapidly declines as the battery discharges, becoming redder and less visible. Regulated LED flashlights may also have user-selectable levels of output appropriate to a task, for example, low light for reading a map and high output for checking a road sign. This would be difficult to do with a single incandescent bulb since efficacy of the lamp drops rapidly at low output. LED flashlights may consume 1 watt or much more from the battery, producing heat as well as light. In contrast to tungsten filaments, which must be hot to produce light, both the light output and the life of an LED decrease with temperature. Heat dissipation for the LED often dictates that small, high-power LED flashlights have aluminium or other high heat-conductivity bodies, reflectors, and other parts to dissipate heat; they can become warm during use. Light output from LED flashlights varies even more widely than for incandescent lights. "Keychain" type lamps operating on button batteries, or lights using a single 5 mm LED, may only produce a few lumens. Even a small LED flashlight operating on an AA cell, but equipped with an LED, can emit 100 lumens. The most powerful LED flashlights produce more than 100,000 lumens and may use multiple LEDs. LEDs are highly efficient at producing colored light compared with incandescent lamps and filters. An LED flashlight may contain different LEDs for white and colored light, selectable by the user for different purposes. Colored LED flashlights are used for signalling, special inspection tasks, forensic examination, or to track the blood trail of wounded game animals. A flashlight may have a red LED intended to preserve dark adaptation of vision. Ultraviolet LEDs may be used for inspection lights, for example, detecting fluorescent dyes added to air conditioning systems to detect leakage, examining paper currency, or checking UV-fluorescing marks on laundry or event ticket holders. Infrared LEDs can be used for illuminators for night-vision systems. LED flashlights may be specified to be compatible with night vision devices. HID A less common type of flashlight uses a high-intensity discharge lamp (HID lamp) as the light source. An HID gas discharge lamp uses a mixture of metal halide salts and noble gas as a filler. HID lamps produce more light than a traditional incandescent flashlight using the same amount of electricity, though not as much as high power LEDs. The lamp lasts longer and is more shock resistant than a regular incandescent bulb, since it lacks the relatively fragile electrical filament found in incandescent bulbs. However, they are much more expensive than incandescent, due to the ballast circuit required to start and operate the lamp. An HID lamp requires a short warm-up time before it reaches full output. LEP LEP stands for Laser Excited Phosphor. The light source is a blue laser diode, which is directed at a phosphor layer to make white light. With the first LEP flashlight available in 2018, there are currently a few dozen LEP flashlights, mainly from China. At the moment, there are 2 types of LEP modules used. The laser light either shines through the phosphor layer to produce white light, or is directed at the layer by a mirror. The mirror-type is built inside a plastic module, while the shine-through models are usually built with a copper/aluminum shell, and much smaller than the plastic type. Accessories Accessories for a flashlight allow the color of the light to be altered or allow light to be dispersed differently. Translucent colored plastic cones slipped over the lens of a flashlight increase the visibility when looking at the side of the light. Such marshalling wands are frequently used for directing automobiles or aircraft at night. Colored lenses placed over the end of the flashlight are used for signalling, for example, in railway yards. Colored light is occasionally useful for hunters tracking wounded game after dusk, or for forensic examination of an area. A red filter helps preserve night vision after the flashlight is turned off, and can be useful to observe animals (such as nesting loggerhead sea turtles) without disturbing them. Detachable light guides, consisting of rigid, bent plastic rods or semirigid or flexible tubes containing optical fibers, are available for some flashlights for inspection inside tanks, or within walls or structures; when not required, the light guide can be removed and the light used for other purposes. Formats and specialized designs A penlight is a small, pen-sized flashlight, often containing two AA or AAA batteries. In some types, the incandescent light bulb has an integral lens that focuses the light, so no reflector is built into the penlight. Others use incandescent bulbs mounted in reflectors. LED penlights are becoming increasingly common. Low-cost units may be disposable with no provision to replace batteries or bulbs and are sometimes imprinted with advertising for promotional purposes. A headlamp is designed to be worn on the head, often having separate lamp and battery components. The battery pack may be attached at the back of the head or in a pocket to improve balance. Headlamps leave the users' hands' free. A headlamp can be clipped to the brim of a hat, or built to mount on a hard hat, instead of using straps; other types resemble eyeglass frames. Similar to the headlamp, an angle-head flashlight emits light perpendicular to the length of the battery tube; it can be clipped to a headband, belt, or webbing or set on a flat surface. Some types allow the user to adjust the angle of the head. The Fulton MX991/U Flashlight was an angle-head flashlight issued to US military personnel; similar style lights remain popular. Tactical lights are sometimes mounted to a handgun or rifle. They allow momentary illumination of a target. They are small enough to be easily rail-mounted to a gun barrel. Tactical lights must withstand the impact of recoil and must be easily controlled while holding the weapon. Although most flashlights are designed for user replacement of the batteries and the bulb as needed, fully sealed disposable flashlights, such as inexpensive keyring lights, are made. When the batteries are depleted or the bulb fails, the entire product is discarded. Diving lamps must be watertight under pressure and are used for night diving and supplemental illumination where surface light cannot reach. The battery compartment of a dive lamp may have a catalyst to recombine any hydrogen gas emitted from the battery since gas cannot be vented in use. People working in hazardous areas with significant concentrations of flammable gases or dusts, such as mines, engine rooms of ships, chemical plants, or grain elevators, use "nonincendive", "intrinsically safe", or "explosion-proof" flashlights constructed so that any spark in the flashlight is not likely to set off an explosion outside the light. The flashlight may require approval by an authority for the particular service and particular gases or dusts expected. The external temperature rise of the flashlight must not exceed the autoignition point of the gas, so substitution of more powerful lamps or batteries may void the approval. Inspection flashlights have permanently mounted light guides containing optical fibers or plastic rods. Another style has a lamp mounted at the end of a flexible cable, or a semirigid or articulated probe. Such lamps are used for inspection inside tanks, or inside structures such as aircraft. Where used for inspecting the interior of tanks containing flammable liquids, the inspection lights may also be rated as flame-proof (explosion-proof) so that they cannot ignite liquids or vapors. Otoscopes and ophthalmoscopes are medical instruments that combine a hand-held light source and magnifying lenses for examination of the ear canal and eyes, respectively. Aboard naval ships, battle lanterns may be used as emergency portable lighting. Installed in major compartments of the ship, a battle lantern can be detached from its mounting and used as portable lighting in the event primary lighting is out of service. Battle lanterns may use either incandescent or LED lamps and may have either disposable primary or rechargeable batteries. Many flashlights are cylindrical in design, with the lamp assembly attached to one end. However, early designs came in a variety of other shapes. Some resembled candlesticks, with a bulb mounted at the top of a battery tube fixed to a flat base, with a handle. Many resembled lanterns, consisting of a battery box with a handle and the lamp and reflector attached to the front. Electric lanterns are used for lighting the broad area immediately around the lantern, as opposed to forming a narrow beam; they can be set down on a level surface or attached to supports. Some electric lanterns use miniature fluorescent lamps for higher efficiency than incandescent bulbs. Portable hand-held electric spotlights can provide larger reflectors and lamps and more powerful batteries than tubular flashlights meant to fit in a pocket. Multifunction portable devices may include a flashlight as one of their features, for example, a portable radio/flashlight combination. Many smartphones have a button or software application available to turn up their screen backlights to full intensity or to switch on the camera flash or video light, providing a "flashlight" function. In addition to utilitarian flashlights, novelty, toy, and ornamental portable electric lights have been made in a myriad of shapes; in the 1890s, one of the earliest portable battery light applications was a type of novelty porcelain tie pin with a concealed bulb and battery. Power sources Batteries The most common power source for flashlights is the battery. Primary battery (disposable) types used in flashlights include button cells, carbon-zinc batteries in both regular and heavy duty types, alkaline, and lithium. Secondary, rechargeable types include lead-acid batteries, NiMH, NiCd batteries and lithium-ion batteries. The choice of batteries plays a determining role in the size, weight, run time, and shape of the flashlight. Flashlight users may prefer a common battery type to simplify replacement. Primary cells are most economical for infrequent use. Some types of lithium primary cell can be stored for years with less risk of leakage compared with zinc-type batteries. Long storage life is useful where flashlights are required only in emergencies. Lithium primary batteries are also useful at lower temperatures than zinc batteries, all of which have water-based electrolytes. Lithium primary batteries have a lower internal resistance than zinc primary batteries, so are more efficient in high-drain flashlights. Flashlights used for extended periods every day may be more economically operated on rechargeable (secondary) batteries. Flashlights designed for rechargeable batteries may allow charging without removing the batteries; for example, a light kept in a vehicle may be trickle-charged and always ready when needed. Some rechargeable flashlights have indicators for the state of charge of the battery. Power-failure lights are designed to keep their batteries charged from a wall plug and to automatically turn on after an AC power failure; the power-failure light can be removed from the wall socket and used as a portable flashlight. Solar powered flashlights use energy from solar cells to charge an on-board battery for later use. Mechanical power One type of mechanically powered flashlight has a winding crank and spring connected to a small electrical generator (dynamo). Some types use the dynamo to charge a capacitor or battery, while others only light while the dynamo is moving. Others generate electricity using electromagnetic induction. They use a strong permanent magnet that can freely slide up and down a tube, passing through a coil of wire as it does. Shaking the flashlight charges a capacitor or a rechargeable battery that supplies current to a light source. Such flashlights can be useful during an emergency, when utility power and batteries may not be available. Dynamo-powered flashlights were popular during the Second World War since replacement batteries were difficult to find. Capacitor At least one manufacturer makes a rechargeable flashlight that uses a supercapacitor to store energy. The capacitor can be recharged more rapidly than a battery and can be recharged many times without loss of capacity; however, the running time is limited by the relative bulk of capacitors compared to electrochemical batteries. Reflectors and lenses A reflector with an approximately parabolic shape concentrates the light emitted by the bulb into a directed beam. Some flashlights allow the user to adjust the relative position of the lamp and reflector, giving a variable-focus effect from a wide floodlight to a narrow beam. Reflectors may be made of polished metal, glass, or plastic with an aluminized reflective finish. Some manufacturers use a pebbled or "orange peel", instead of a smooth reflector, to improve the uniformity of the light beam emitted. Where multiple LEDs are used, each one may be put in its own parabolic reflector. Flashlights using a "total internal reflection" assembly have a transparent optical element (light pipe) to guide light from the source into a beam; no reflector surface is required. For a given size of light source, a larger reflector or lens allows a tighter beam to be produced, while capturing the same fraction of the emitted light. Some flashlights use Fresnel lenses, which allow the weight of the lens to be reduced. The reflector may have a flat transparent cover to keep out dirt and moisture, but some designs have a plastic or glass "bulls-eye" lens to form a concentrated beam. The lens or reflector cover must resist impacts and the heat of the lamp, and must not lose too much of the transmitted light to reflection or absorption. Very small flashlights may not have a reflector or lens separate from the lamp. Some types of penlight bulbs or small LEDs have a built-in lens. A reflector forms a narrow beam called the "throw" in hobbyist parlance, while light emitted forward misses the reflector and forms a wide flood or "spill" of light. Because LEDs emit most light in a hemisphere, lens lights with the LED facing forward or reflector lights with it facing backwards radiate less spill. Variable focus "zoom" or "flood to throw" lights may move the reflector or lens or they may move the emitter; moving the emitter presents the designer with the problem of maintaining heat dissipation for the LED. Control switch The original 1890s flashlights used a metal ring around the fiber body of the flashlight as one contact of a switch; the second contact was a movable metal loop that could be flipped down to touch the ring, completing the circuit. A wide variety of mechanical switch designs using slide switches, rocker switches, or side-mounted or end-mounted pushbuttons has been used in flashlights. A common combination is a slide switch that allows the light to be left on for an extended time, combined with a momentary button for intermittent use or signalling. (On earlier models, the button was a switch and the slider simply locked the button down.) Since voltages and currents are low, switch design is limited only by the available space and desired cost of production. Switches may be covered with a flexible rubber boot to exclude dirt and moisture and may be backlit for easy location. Another common type of switch relies on twisting the head of the light. Weapon-mounted lights may have remote switches for convenience in operation. Electronic controls allow the user to select variable output levels or different operating modes such as pre-programmed flashing beacon or strobe modes. Electronic controls may be operated by buttons, sliders, magnets, rotating heads, or rotating control rings. Some models of flashlight include an acceleration sensor to allow them to respond to shaking, or to select modes based on what direction the light is held when switched on. At least one manufacturer allows user programming of the features of the flashlight through a USB port. An electronic control may also provide an indication of remaining battery capacity, voltage, or provide information regarding recharging or automatic step-down of brightness as the battery nears full discharge. Materials Early flashlights used vulcanized fiber or hard rubber tubes with metal end caps. Many other materials including drawn steel, plated brass, copper, or silver, and even wood and leather have been used. Modern flashlights are generally made of plastic or aluminum. Plastics range from low-cost polystyrene and polyethylene to more complex mixtures of ABS or glass-reinforced epoxies. Some manufacturers have proprietary plastic formulations for their products. A desirable plastic for manufacturing flashlights allows for ease of molding and adequate mechanical properties of the finished flashlight case. Aluminum, either plain, painted or anodized, is a popular choice. It is electrically conductive, can be easily machined, and dissipates heat well. Several standard alloys of aluminum are used. Other metals include copper, stainless steel, and titanium, which can be polished to provide a decorative finish. Zinc can be die-cast into intricate shapes. Magnesium and its alloys provide strength and heat dissipation similar to aluminum with less weight, but they corrode easily. Metals may be drawn into a tubular shape, or tubular extruded stock can be machined to add threads for the head and tail cap, knurling for grip, and decorative and functional flats or holes in the body. LED flashlights may have cooling fins machined into their metal cases. Plastics are often injection molded into nearly final shape, requiring only a few more process steps to complete assembly. Metal cases provide better heat dissipation for the LED, but plastics are not electrically conductive and may resist corrosion and wear. Ratings and standards Safety regulations Industrial, marine, public safety, and military organizations develop specifications for flashlights in specialized roles. Typically, light output, overall dimensions, and battery compatibility and durability are required to meet minimum limits. Flashlights may be tested for impact resistance, water and chemical resistance, and the lifespan of the control switch. Flashlights intended for use in hazardous areas with flammable gas or dust are tested to ensure they cannot set off an explosion. Flashlights approved for flammable gas areas have markings indicating the approving agency (MSHA, ATEX, UL, etc.) and symbols for the conditions that were tested. Flashlights for hazardous areas may be designed to automatically disconnect the lamp if the bulb is broken, to prevent ignition of flammable gas. Regulations for ships and aircraft specify the number and general properties of flashlights included as part of the standard safety equipment of the vessel. Flashlights for small boats may be required to be waterproof and to float. Uniformed services may issue particular models of flashlights, or may provide minimum performance standards for their members to follow when purchasing their own flashlights. Performance standards The United States Army former standard MIL-F-3747E described the performance standard for plastic flashlights using two or three D-cell dry batteries, in either straight or angle form, and in standard, explosion-proof, heat-resistant, traffic direction, and inspection types. The standard described only incandescent lamp flashlights and was withdrawn in 1996. In the United States, ANSI in 2009 published FL1 Flashlight basic performance standard. This voluntary standard defines test procedures and conditions for total light output, beam intensity, working distance, impact and water resistance, and battery running time to 10% of initial light output. The FL1 standard gives definitions for terms used in marketing flashlights, with the intention of allowing the consumer to compare products tested to the standard. The standard recommends particular graphic symbols and wording for the product package, so that the consumer can identify products tested to the standard. Testing may be carried out by the manufacturer itself or by a third-party test laboratory. The FL1 standard requires measurements reported on the packaging to be made with the type of batteries packaged with the flashlight, or with an identified type of battery. Initial light output is measured with an integrating sphere photometer, 30 seconds after the light is switched on with fresh (or newly charged) batteries. The total light emitted is reported in lumens. Luminous intensity is determined by measuring the brightest spot in the beam produced by the flashlight, in candelas. Since this is a measure of all the light emitted in a solid angle (the "cone" of light in a particular direction), the beam intensity is independent of distance. The working distance is defined as the distance at which the maximum light falling on a surface (illuminance) would fall to 0.25 lux. This is comparable to a full moon on a clear night. The distance is calculated from the square root of (the beam intensity in candelas divided by 0.25 lux); for example, a beam intensity of 1000 candelas produces a working range rating of the square root of (1000/0.25), or 63 meters. The result is reported in meters or feet. The working distance is from the point of view of the user of the flashlight. A light directly pointed at an observer may be visible against a dark background for many times this distance, especially if the observer has night-vision equipment. Run time is measured using the supplied or specified batteries and letting the light run until the intensity of the beam has dropped to 10% of the value 30 seconds after switching on. The standard does not evaluate the behavior of the flashlight output during run time. A regulated flashlight may run at only a slowly declining output and then abruptly cut off, but unregulated types may have steeply-declining light output after only a short time. Manufacturers of headlamps may use a different standard which rates run times until light output falls to 1 lux at 2 meters distance; this value is not comparable to the FL 1 runtime measurement. Impact resistance is measured by dropping the flashlight in six different orientations and observing that it still functions and has no large cracks or breaks in it; the height used in the test is reported. Water resistance, if specified, is evaluated after impact testing; no water is to be visible inside the unit and it must remain functional. Ratings are given in IP Code terms, where jet spray corresponds to IP X6, brief immersion to IPX7, 30 minutes immersion at 1 meter or more is IP X8; (the depth is reported if greater than 1 meter). An IP X8 rating by FL1 does not imply that the lamp is suitable for use as a diver's light since the test protocol examines function of the light only after immersion, not during immersion. The consumer must decide how well the ANSI test conditions match their requirements, but all manufacturers testing to the FL1 standard can be compared on a uniform basis. The light measurements are more directly related to the use of flashlights than is the nominal power input to the lamp (watts), since different LED and incandescent lamp types vary widely in the amount of light produced per watt. Even the same LED or lamp in different optical systems will show different beam characteristics. The visibility of objects depends on many factors as well as the amount of light emitted by the flashlight. ANSI standard FL1 does not specify measurements of the beam width angle but the candela intensity and total lumen ratings can be used by the consumer to assess the beam characteristics. Where two flashlights have similar total light (lumen) measures, the unit with the higher candela rating produces a more concentrated beam of light, suitable for lighting distant objects; it will also have a higher working distance. If two lights have similar candela ratings, the light with higher lumen value will produce a wider beam and will light a wider area overall. A beam width (containing most of the power of the beam, or "hot spot") of a few degrees corresponds to a spot light, useful for searching for distant objects; beam widths of 20 degrees or more are described as flood lights, suitable for lighting a wide nearby area. Typically even a flashlight beam with a small hot spot will have some light visible as "spill" around the spot. In 2018, in the United States, Underwriter's Laboratories published UL standard 1576 for flashlights and lanterns, outlining safety requirements and performance tests. Applications	Technology	Lighting	null
243840	https://en.wikipedia.org/wiki/Curlew	Curlew	The curlews () are a group of nine species of birds in the genus Numenius, characterised by their long, slender, downcurved bills and mottled brown plumage. The English name is imitative of the Eurasian curlew's call, but may have been influenced by the Old French corliu, "messenger", from courir , "to run". It was first recorded in 1377 in Langland's Piers Plowman "Fissch to lyue in þe flode..Þe corlue by kynde of þe eyre". In Europe, "curlew" usually refers to one species, the Eurasian curlew (Numenius arquata). Description They are one of the most ancient lineages of scolopacid waders, together with the godwits which look similar but have straight bills. Curlews feed on mud or very soft ground, searching for worms and other invertebrates with their long bills. They will also take crabs and similar items. Distribution Curlews enjoy a worldwide distribution. Most species exhibit strong migratory habits and consequently one or more species can be encountered at different times of the year in Europe, Ireland, Britain, Iberia, Iceland, Africa, Southeast Asia, Siberia, North America, South America and Australasia. The distribution of curlews has altered considerably in the past hundred years as a result of changing agricultural practices. For instance, Eurasian curlew populations have suffered due to draining of marshes for farmland, whereas long-billed curlews have shown an increase in breeding densities around areas grazed by livestock. , there were only a small number of Eurasian curlews still breeding in Ireland, raising concerns that the bird will become extinct in that country. The stone-curlews are not true curlews (family Scolopacidae) but members of the family Burhinidae, which is in the same order Charadriiformes, but only distantly related within that. Taxonomy The genus Numenius was erected by the French scientist Mathurin Jacques Brisson in his Ornithologie published in 1760. The type species is the Eurasian curlew (Numenius arquata). The Swedish naturalist Carl Linnaeus had introduced the genus Numenius in the 6th edition of his Systema Naturae published in 1748, but Linnaeus dropped the genus in the important tenth edition of 1758 and put the curlews together with the woodcocks in the genus Scolopax. As the publication date of Linnaeus's sixth edition was before the 1758 starting point of the International Commission on Zoological Nomenclature, Brisson and not Linnaeus is considered as the authority for the genus. The name Numenius is from Ancient Greek noumenios, a bird mentioned by Hesychius. It is associated with the curlews because it appears to be derived from neos, "new" and mene "moon", referring to the crescent-shaped bill. The genus contains nine species: The following cladogram shows the genetic relationships between the species. It is based on a study published in 2023. The Late Eocene (Montmartre Formation, some 35 mya) fossil Limosa gypsorum of France was originally placed in Numenius and may in fact belong there. Apart from that, a Late Pleistocene curlew from San Josecito Cave, Mexico has been described. This fossil was initially placed in a distinct genus, Palnumenius, but was actually a chronospecies or paleosubspecies related to the long-billed curlew. The upland sandpiper (Bartramia longicauda) is an odd bird which is the closest relative of the curlews. It is distinguished from them by its yellow legs, long tail, and shorter, less curved bill.	Biology and health sciences	Charadriiformes	Animals
243844	https://en.wikipedia.org/wiki/Extraterrestrial%20intelligence	Extraterrestrial intelligence	Extraterrestrial intelligence (ETI) refers to hypothetical intelligent extraterrestrial life. No such life has ever been verifiably observed to exist. The question of whether other inhabited worlds might exist has been debated since ancient times. The modern form of the concept emerged when the Copernican Revolution demonstrated that the Earth was a planet revolving around the Sun, and other planets were, conversely, other worlds. The question of whether other inhabited planets or moons exist was a natural consequence of this new understanding. It has become one of the most speculative questions in science and is a central theme of science fiction and popular culture. An alternative name for it is "Extraterrestrial Technological Instantiations" (ETI). The term was coined to avoid the use of terms such as "civilizations" "species" and "intelligence", as those may prove to be ambiguous and open to interpretation, or simply inapplicable in its local context. Intelligence Intelligence is, along with the more precise concept of sapience, used to describe extraterrestrial life with similar cognitive abilities as humans. Another interchangeable term is sophoncy, being wise or wiser, first coined by Karen Anderson and published in the 1966 works by her husband Poul Anderson. Sentience, like consciousness, is a concept sometimes mistakenly used to refer to the concept of intelligence and sapience, since it does not exclude forms of life that are non-sapient (or more broadly non-intelligent or non-conscious). The term extraterrestrial civilization frames a more particular case of extraterrestrial intelligence. It is the possible long-term result of intelligent and specifically sapient extraterrestrial life. Probability The Copernican principle is generalized to the relativistic concept that humans are not privileged observers of the universe. Many prominent scientists, including Stephen Hawking have proposed that the sheer scale of the universe makes it improbable for intelligent life not to have emerged elsewhere. However, Fermi's Paradox highlights the apparent contradiction between high estimates of the probability of the existence of extraterrestrial civilization and humanity's lack of contact with, or evidence for, such civilizations. So far, there is no observation of extraterrestrial life, including intelligent extraterrestrial life. The Kardashev scale is a speculative method of measuring a civilization's level of technological advancement, based on the amount of energy a civilization is able to utilize. The Drake equation is a probabilistic framework used to estimate the number of active, communicative extraterrestrial civilizations in the Milky Way galaxy. Search for extraterrestrial intelligence There has been a search for signals from extraterrestrial intelligence for several decades, with no significant results. Active SETI (Active Search for Extra-Terrestrial Intelligence) is the attempt to send messages to intelligent extraterrestrial life. Active SETI messages are usually sent in the form of radio signals. Physical messages like that of the Pioneer plaque may also be considered an active SETI message. Communication with extraterrestrial intelligence (CETI) is a branch of the search for extraterrestrial intelligence that focuses on composing and deciphering messages that could theoretically be understood by another technological civilization. The best-known CETI experiment was the 1974 Arecibo message composed by Frank Drake and Carl Sagan. There are multiple independent organizations and individuals engaged in CETI research. The U.S. government's position, in line with that of most relevant experts, is that "chances of contact with an extraterrestrial intelligence are extremely small, given the distances involved." This line of thinking has led some to conclude that first contact might be made with extraterrestrial artificial intelligence, rather than with biological beings. The Wow! signal remains the best candidate for an extraterrestrial radio signal ever detected, though the fact that no similar signal has ever been observed again makes attribution of the signal to any cause difficult if not impossible. On 14 June 2022 astronomers working with China's FAST telescope reported the possibility of having detected artificial (presumably alien) signals, but cautions that further studies are required to determine if some kind of natural radio interference may be the source. On 18 June 2022 Dan Werthimer, chief scientist for several SETI-related projects, reportedly noted that “These signals are from radio interference; they are due to radio pollution from earthlings, not from E.T.” Potential cultural impact of extraterrestrial contact The potential changes from extraterrestrial contact could vary greatly in magnitude and type, based on the extraterrestrial civilization's level of technological advancement, degree of benevolence or malevolence, and level of mutual comprehension between itself and humanity. Some theories suggest that an extraterrestrial civilization could be advanced enough to dispense with biology, living instead inside of advanced computers. The medium through which humanity is contacted, be it electromagnetic radiation, direct physical interaction, extraterrestrial artefact, or otherwise, may also influence the results of contact. Incorporating these factors, various systems have been created to assess the implications of extraterrestrial contact. The implications of extraterrestrial contact, particularly with a technologically superior civilization, have often been likened to the meeting of two vastly different human cultures on Earth, a historical precedent being the Columbian Exchange. Such meetings have generally led to the destruction of the civilization receiving contact (as opposed to the "contactor", which initiates contact), and therefore destruction of human civilization is a possible outcome. However, the absence of any such contact to date means such conjecture is largely speculative. UFOlogy The extraterrestrial hypothesis is the idea that some UFOs are vehicles containing or sent by extraterrestrial beings (usually called aliens in this context). As an explanation for UFOs, ETI is sometimes contrasted with EDI (extradimensional intelligence), for example by J. Allen Hynek. In 2023, House lawmakers held a hearing to examine how the executive branch handles reports of UFOs. In culture The theories and reception of the probability of intelligent life has been a recurring cultural element, particularly of popular culture since the prospect and achievement of spaceflight. New Mexico has even declared in 2003 the 14th of February as the Extraterrestrial Culture Day.	Physical sciences	Astronomy basics	Astronomy
243849	https://en.wikipedia.org/wiki/Projective%20geometry	Projective geometry	In mathematics, projective geometry is the study of geometric properties that are invariant with respect to projective transformations. This means that, compared to elementary Euclidean geometry, projective geometry has a different setting (projective space) and a selective set of basic geometric concepts. The basic intuitions are that projective space has more points than Euclidean space, for a given dimension, and that geometric transformations are permitted that transform the extra points (called "points at infinity") to Euclidean points, and vice versa. Properties meaningful for projective geometry are respected by this new idea of transformation, which is more radical in its effects than can be expressed by a transformation matrix and translations (the affine transformations). The first issue for geometers is what kind of geometry is adequate for a novel situation. Unlike in Euclidean geometry, the concept of an angle does not apply in projective geometry, because no measure of angles is invariant with respect to projective transformations, as is seen in perspective drawing from a changing perspective. One source for projective geometry was indeed the theory of perspective. Another difference from elementary geometry is the way in which parallel lines can be said to meet in a point at infinity, once the concept is translated into projective geometry's terms. Again this notion has an intuitive basis, such as railway tracks meeting at the horizon in a perspective drawing. See Projective plane for the basics of projective geometry in two dimensions. While the ideas were available earlier, projective geometry was mainly a development of the 19th century. This included the theory of complex projective space, the coordinates used (homogeneous coordinates) being complex numbers. Several major types of more abstract mathematics (including invariant theory, the Italian school of algebraic geometry, and Felix Klein's Erlangen programme resulting in the study of the classical groups) were motivated by projective geometry. It was also a subject with many practitioners for its own sake, as synthetic geometry. Another topic that developed from axiomatic studies of projective geometry is finite geometry. The topic of projective geometry is itself now divided into many research subtopics, two examples of which are projective algebraic geometry (the study of projective varieties) and projective differential geometry (the study of differential invariants of the projective transformations). Overview Projective geometry is an elementary non-metrical form of geometry, meaning that it does not support any concept of distance. In two dimensions it begins with the study of configurations of points and lines. That there is indeed some geometric interest in this sparse setting was first established by Desargues and others in their exploration of the principles of perspective art. In higher dimensional spaces there are considered hyperplanes (that always meet), and other linear subspaces, which exhibit the principle of duality. The simplest illustration of duality is in the projective plane, where the statements "two distinct points determine a unique line" (i.e. the line through them) and "two distinct lines determine a unique point" (i.e. their point of intersection) show the same structure as propositions. Projective geometry can also be seen as a geometry of constructions with a straight-edge alone, excluding compass constructions, common in straightedge and compass constructions. As such, there are no circles, no angles, no measurements, no parallels, and no concept of intermediacy (or "betweenness"). It was realised that the theorems that do apply to projective geometry are simpler statements. For example, the different conic sections are all equivalent in (complex) projective geometry, and some theorems about circles can be considered as special cases of these general theorems. During the early 19th century the work of Jean-Victor Poncelet, Lazare Carnot and others established projective geometry as an independent field of mathematics. Its rigorous foundations were addressed by Karl von Staudt and perfected by Italians Giuseppe Peano, Mario Pieri, Alessandro Padoa and Gino Fano during the late 19th century. Projective geometry, like affine and Euclidean geometry, can also be developed from the Erlangen program of Felix Klein; projective geometry is characterized by invariants under transformations of the projective group. After much work on the very large number of theorems in the subject, therefore, the basics of projective geometry became understood. The incidence structure and the cross-ratio are fundamental invariants under projective transformations. Projective geometry can be modeled by the affine plane (or affine space) plus a line (hyperplane) "at infinity" and then treating that line (or hyperplane) as "ordinary". An algebraic model for doing projective geometry in the style of analytic geometry is given by homogeneous coordinates. On the other hand, axiomatic studies revealed the existence of non-Desarguesian planes, examples to show that the axioms of incidence can be modelled (in two dimensions only) by structures not accessible to reasoning through homogeneous coordinate systems. In a foundational sense, projective geometry and ordered geometry are elementary since they each involve a minimal set of axioms and either can be used as the foundation for affine and Euclidean geometry. Projective geometry is not "ordered" and so it is a distinct foundation for geometry. Description Projective geometry is less restrictive than either Euclidean geometry or affine geometry. It is an intrinsically non-metrical geometry, meaning that facts are independent of any metric structure. Under the projective transformations, the incidence structure and the relation of projective harmonic conjugates are preserved. A projective range is the one-dimensional foundation. Projective geometry formalizes one of the central principles of perspective art: that parallel lines meet at infinity, and therefore are drawn that way. In essence, a projective geometry may be thought of as an extension of Euclidean geometry in which the "direction" of each line is subsumed within the line as an extra "point", and in which a "horizon" of directions corresponding to coplanar lines is regarded as a "line". Thus, two parallel lines meet on a horizon line by virtue of their incorporating the same direction. Idealized directions are referred to as points at infinity, while idealized horizons are referred to as lines at infinity. In turn, all these lines lie in the plane at infinity. However, infinity is a metric concept, so a purely projective geometry does not single out any points, lines or planes in this regard—those at infinity are treated just like any others. Because a Euclidean geometry is contained within a projective geometry—with projective geometry having a simpler foundation—general results in Euclidean geometry may be derived in a more transparent manner, where separate but similar theorems of Euclidean geometry may be handled collectively within the framework of projective geometry. For example, parallel and nonparallel lines need not be treated as separate cases; rather an arbitrary projective plane is singled out as the ideal plane and located "at infinity" using homogeneous coordinates. Additional properties of fundamental importance include Desargues' Theorem and the Theorem of Pappus. In projective spaces of dimension 3 or greater there is a construction that allows one to prove Desargues' Theorem. But for dimension 2, it must be separately postulated. Using Desargues' Theorem, combined with the other axioms, it is possible to define the basic operations of arithmetic, geometrically. The resulting operations satisfy the axioms of a field – except that the commutativity of multiplication requires Pappus's hexagon theorem. As a result, the points of each line are in one-to-one correspondence with a given field, , supplemented by an additional element, ∞, such that , , , , , , except that , , , , and remain undefined. Projective geometry also includes a full theory of conic sections, a subject also extensively developed in Euclidean geometry. There are advantages to being able to think of a hyperbola and an ellipse as distinguished only by the way the hyperbola lies across the line at infinity; and that a parabola is distinguished only by being tangent to the same line. The whole family of circles can be considered as conics passing through two given points on the line at infinity — at the cost of requiring complex coordinates. Since coordinates are not "synthetic", one replaces them by fixing a line and two points on it, and considering the linear system of all conics passing through those points as the basic object of study. This method proved very attractive to talented geometers, and the topic was studied thoroughly. An example of this method is the multi-volume treatise by H. F. Baker. History The first geometrical properties of a projective nature were discovered during the 3rd century by Pappus of Alexandria. Filippo Brunelleschi (1404–1472) started investigating the geometry of perspective during 1425 (see for a more thorough discussion of the work in the fine arts that motivated much of the development of projective geometry). Johannes Kepler (1571–1630) and Girard Desargues (1591–1661) independently developed the concept of the "point at infinity". Desargues developed an alternative way of constructing perspective drawings by generalizing the use of vanishing points to include the case when these are infinitely far away. He made Euclidean geometry, where parallel lines are truly parallel, into a special case of an all-encompassing geometric system. Desargues's study on conic sections drew the attention of 16-year-old Blaise Pascal and helped him formulate Pascal's theorem. The works of Gaspard Monge at the end of 18th and beginning of 19th century were important for the subsequent development of projective geometry. The work of Desargues was ignored until Michel Chasles chanced upon a handwritten copy during 1845. Meanwhile, Jean-Victor Poncelet had published the foundational treatise on projective geometry during 1822. Poncelet examined the projective properties of objects (those invariant under central projection) and, by basing his theory on the concrete pole and polar relation with respect to a circle, established a relationship between metric and projective properties. The non-Euclidean geometries discovered soon thereafter were eventually demonstrated to have models, such as the Klein model of hyperbolic space, relating to projective geometry. In 1855 A. F. Möbius wrote an article about permutations, now called Möbius transformations, of generalised circles in the complex plane. These transformations represent projectivities of the complex projective line. In the study of lines in space, Julius Plücker used homogeneous coordinates in his description, and the set of lines was viewed on the Klein quadric, one of the early contributions of projective geometry to a new field called algebraic geometry, an offshoot of analytic geometry with projective ideas. Projective geometry was instrumental in the validation of speculations of Lobachevski and Bolyai concerning hyperbolic geometry by providing models for the hyperbolic plane: for example, the Poincaré disc model where generalised circles perpendicular to the unit circle correspond to "hyperbolic lines" (geodesics), and the "translations" of this model are described by Möbius transformations that map the unit disc to itself. The distance between points is given by a Cayley–Klein metric, known to be invariant under the translations since it depends on cross-ratio, a key projective invariant. The translations are described variously as isometries in metric space theory, as linear fractional transformations formally, and as projective linear transformations of the projective linear group, in this case . The work of Poncelet, Jakob Steiner and others was not intended to extend analytic geometry. Techniques were supposed to be synthetic: in effect projective space as now understood was to be introduced axiomatically. As a result, reformulating early work in projective geometry so that it satisfies current standards of rigor can be somewhat difficult. Even in the case of the projective plane alone, the axiomatic approach can result in models not describable via linear algebra. This period in geometry was overtaken by research on the general algebraic curve by Clebsch, Riemann, Max Noether and others, which stretched existing techniques, and then by invariant theory. Towards the end of the century, the Italian school of algebraic geometry (Enriques, Segre, Severi) broke out of the traditional subject matter into an area demanding deeper techniques. During the later part of the 19th century, the detailed study of projective geometry became less fashionable, although the literature is voluminous. Some important work was done in enumerative geometry in particular, by Schubert, that is now considered as anticipating the theory of Chern classes, taken as representing the algebraic topology of Grassmannians. Projective geometry later proved key to Paul Dirac's invention of quantum mechanics. At a foundational level, the discovery that quantum measurements could fail to commute had disturbed and dissuaded Heisenberg, but past study of projective planes over noncommutative rings had likely desensitized Dirac. In more advanced work, Dirac used extensive drawings in projective geometry to understand the intuitive meaning of his equations, before writing up his work in an exclusively algebraic formalism. Classification There are many projective geometries, which may be divided into discrete and continuous: a discrete geometry comprises a set of points, which may or may not be finite in number, while a continuous geometry has infinitely many points with no gaps in between. The only projective geometry of dimension 0 is a single point. A projective geometry of dimension 1 consists of a single line containing at least 3 points. The geometric construction of arithmetic operations cannot be performed in either of these cases. For dimension 2, there is a rich structure in virtue of the absence of Desargues' Theorem. The smallest 2-dimensional projective geometry (that with the fewest points) is the Fano plane, which has 3 points on every line, with 7 points and 7 lines in all, having the following collinearities: [ABC] [ADE] [AFG] [BDG] [BEF] [CDF] [CEG] with homogeneous coordinates , , , , , , , or, in affine coordinates, , , , , , and . The affine coordinates in a Desarguesian plane for the points designated to be the points at infinity (in this example: C, E and G) can be defined in several other ways. In standard notation, a finite projective geometry is written where: is the projective (or geometric) dimension, and is one less than the number of points on a line (called the order of the geometry). Thus, the example having only 7 points is written . The term "projective geometry" is used sometimes to indicate the generalised underlying abstract geometry, and sometimes to indicate a particular geometry of wide interest, such as the metric geometry of flat space which we analyse through the use of homogeneous coordinates, and in which Euclidean geometry may be embedded (hence its name, Extended Euclidean plane). The fundamental property that singles out all projective geometries is the elliptic incidence property that any two distinct lines and in the projective plane intersect at exactly one point . The special case in analytic geometry of parallel lines is subsumed in the smoother form of a line at infinity on which lies. The line at infinity is thus a line like any other in the theory: it is in no way special or distinguished. (In the later spirit of the Erlangen programme one could point to the way the group of transformations can move any line to the line at infinity). The parallel properties of elliptic, Euclidean and hyperbolic geometries contrast as follows: Given a line and a point not on the line, Elliptic there exists no line through that does not meet Euclidean there exists exactly one line through that does not meet Hyperbolic there exists more than one line through that does not meet The parallel property of elliptic geometry is the key idea that leads to the principle of projective duality, possibly the most important property that all projective geometries have in common. Duality In 1825, Joseph Gergonne noted the principle of duality characterizing projective plane geometry: given any theorem or definition of that geometry, substituting point for line, lie on for pass through, collinear for concurrent, intersection for join, or vice versa, results in another theorem or valid definition, the "dual" of the first. Similarly in 3 dimensions, the duality relation holds between points and planes, allowing any theorem to be transformed by swapping point and plane, is contained by and contains. More generally, for projective spaces of dimension N, there is a duality between the subspaces of dimension R and dimension . For , this specializes to the most commonly known form of duality—that between points and lines. The duality principle was also discovered independently by Jean-Victor Poncelet. To establish duality only requires establishing theorems which are the dual versions of the axioms for the dimension in question. Thus, for 3-dimensional spaces, one needs to show that (1) every point lies in 3 distinct planes, (2) every two planes intersect in a unique line and a dual version of (3*) to the effect: if the intersection of plane P and Q is coplanar with the intersection of plane R and S, then so are the respective intersections of planes P and R, Q and S (assuming planes P and S are distinct from Q and R). In practice, the principle of duality allows us to set up a dual correspondence between two geometric constructions. The most famous of these is the polarity or reciprocity of two figures in a conic curve (in 2 dimensions) or a quadric surface (in 3 dimensions). A commonplace example is found in the reciprocation of a symmetrical polyhedron in a concentric sphere to obtain the dual polyhedron. Another example is Brianchon's theorem, the dual of the already mentioned Pascal's theorem, and one of whose proofs simply consists of applying the principle of duality to Pascal's. Here are comparative statements of these two theorems (in both cases within the framework of the projective plane): Pascal: If all six vertices of a hexagon lie on a conic, then the intersections of its opposite sides (regarded as full lines, since in the projective plane there is no such thing as a "line segment") are three collinear points. The line joining them is then called the Pascal line of the hexagon. Brianchon: If all six sides of a hexagon are tangent to a conic, then its diagonals (i.e. the lines joining opposite vertices) are three concurrent lines. Their point of intersection is then called the Brianchon point of the hexagon. (If the conic degenerates into two straight lines, Pascal's becomes Pappus's theorem, which has no interesting dual, since the Brianchon point trivially becomes the two lines' intersection point.) Axioms of projective geometry Any given geometry may be deduced from an appropriate set of axioms. Projective geometries are characterised by the "elliptic parallel" axiom, that any two planes always meet in just one line, or in the plane, any two lines always meet in just one point. In other words, there are no such things as parallel lines or planes in projective geometry. Many alternative sets of axioms for projective geometry have been proposed (see for example Coxeter 2003, Hilbert & Cohn-Vossen 1999, Greenberg 1980). Whitehead's axioms These axioms are based on Whitehead, "The Axioms of Projective Geometry". There are two types, points and lines, and one "incidence" relation between points and lines. The three axioms are: G1: Every line contains at least 3 points G2: Every two distinct points, A and B, lie on a unique line, AB. G3: If lines AB and CD intersect, then so do lines AC and BD (where it is assumed that A and D are distinct from B and C). The reason each line is assumed to contain at least 3 points is to eliminate some degenerate cases. The spaces satisfying these three axioms either have at most one line, or are projective spaces of some dimension over a division ring, or are non-Desarguesian planes. Additional axioms One can add further axioms restricting the dimension or the coordinate ring. For example, Coxeter's Projective Geometry, references Veblen in the three axioms above, together with a further 5 axioms that make the dimension 3 and the coordinate ring a commutative field of characteristic not 2. Axioms using a ternary relation One can pursue axiomatization by postulating a ternary relation, [ABC] to denote when three points (not all necessarily distinct) are collinear. An axiomatization may be written down in terms of this relation as well: C0: [ABA] C1: If A and B are distinct points such that [ABC] and [ABD] then [BDC] C2: If A and B are distinct points then there exists a third distinct point C such that [ABC] C3: If A and C are distinct points, and B and D are distinct points, with [BCE] and [ADE] but not [ABE], then there is a point F such that [ACF] and [BDF]. For two distinct points, A and B, the line AB is defined as consisting of all points C for which [ABC]. The axioms C0 and C1 then provide a formalization of G2; C2 for G1 and C3 for G3. The concept of line generalizes to planes and higher-dimensional subspaces. A subspace, AB...XY may thus be recursively defined in terms of the subspace AB...X as that containing all the points of all lines YZ, as Z ranges over AB...X. Collinearity then generalizes to the relation of "independence". A set of points is independent, [AB...Z] if is a minimal generating subset for the subspace AB...Z. The projective axioms may be supplemented by further axioms postulating limits on the dimension of the space. The minimum dimension is determined by the existence of an independent set of the required size. For the lowest dimensions, the relevant conditions may be stated in equivalent form as follows. A projective space is of: (L1) at least dimension 0 if it has at least 1 point, (L2) at least dimension 1 if it has at least 2 distinct points (and therefore a line), (L3) at least dimension 2 if it has at least 3 non-collinear points (or two lines, or a line and a point not on the line), (L4) at least dimension 3 if it has at least 4 non-coplanar points. The maximum dimension may also be determined in a similar fashion. For the lowest dimensions, they take on the following forms. A projective space is of: (M1) at most dimension 0 if it has no more than 1 point, (M2) at most dimension 1 if it has no more than 1 line, (M3) at most dimension 2 if it has no more than 1 plane, and so on. It is a general theorem (a consequence of axiom (3)) that all coplanar lines intersect—the very principle that projective geometry was originally intended to embody. Therefore, property (M3) may be equivalently stated that all lines intersect one another. It is generally assumed that projective spaces are of at least dimension 2. In some cases, if the focus is on projective planes, a variant of M3 may be postulated. The axioms of (Eves 1997: 111), for instance, include (1), (2), (L3) and (M3). Axiom (3) becomes vacuously true under (M3) and is therefore not needed in this context. Axioms for projective planes In incidence geometry, most authors give a treatment that embraces the Fano plane as the smallest finite projective plane. An axiom system that achieves this is as follows: (P1) Any two distinct points lie on a line that is unique. (P2) Any two distinct lines meet at a point that is unique. (P3) There exist at least four points of which no three are collinear. Coxeter's Introduction to Geometry gives a list of five axioms for a more restrictive concept of a projective plane that is attributed to Bachmann, adding Pappus's theorem to the list of axioms above (which eliminates non-Desarguesian planes) and excluding projective planes over fields of characteristic 2 (those that do not satisfy Fano's axiom). The restricted planes given in this manner more closely resemble the real projective plane. Perspectivity and projectivity Given three non-collinear points, there are three lines connecting them, but with four points, no three collinear, there are six connecting lines and three additional "diagonal points" determined by their intersections. The science of projective geometry captures this surplus determined by four points through a quaternary relation and the projectivities which preserve the complete quadrangle configuration. An harmonic quadruple of points on a line occurs when there is a complete quadrangle two of whose diagonal points are in the first and third position of the quadruple, and the other two positions are points on the lines joining two quadrangle points through the third diagonal point. A spatial perspectivity of a projective configuration in one plane yields such a configuration in another, and this applies to the configuration of the complete quadrangle. Thus harmonic quadruples are preserved by perspectivity. If one perspectivity follows another the configurations follow along. The composition of two perspectivities is no longer a perspectivity, but a projectivity. While corresponding points of a perspectivity all converge at a point, this convergence is not true for a projectivity that is not a perspectivity. In projective geometry the intersection of lines formed by corresponding points of a projectivity in a plane are of particular interest. The set of such intersections is called a projective conic, and in acknowledgement of the work of Jakob Steiner, it is referred to as a Steiner conic. Suppose a projectivity is formed by two perspectivities centered on points A and B, relating x to X by an intermediary p: The projectivity is then Then given the projectivity the induced conic is Given a conic C and a point P not on it, two distinct secant lines through P intersect C in four points. These four points determine a quadrangle of which P is a diagonal point. The line through the other two diagonal points is called the polar of P and P is the pole of this line. Alternatively, the polar line of P is the set of projective harmonic conjugates of P on a variable secant line passing through P and C.	Mathematics	Non-Euclidean geometry	null
243890	https://en.wikipedia.org/wiki/Affine%20geometry	Affine geometry	In mathematics, affine geometry is what remains of Euclidean geometry when ignoring (mathematicians often say "forgetting") the metric notions of distance and angle. As the notion of parallel lines is one of the main properties that is independent of any metric, affine geometry is often considered as the study of parallel lines. Therefore, Playfair's axiom (Given a line and a point not on , there is exactly one line parallel to that passes through .) is fundamental in affine geometry. Comparisons of figures in affine geometry are made with affine transformations, which are mappings that preserve alignment of points and parallelism of lines. Affine geometry can be developed in two ways that are essentially equivalent. In synthetic geometry, an affine space is a set of points to which is associated a set of lines, which satisfy some axioms (such as Playfair's axiom). Affine geometry can also be developed on the basis of linear algebra. In this context an affine space is a set of points equipped with a set of transformations (that is bijective mappings), the translations, which forms a vector space (over a given field, commonly the real numbers), and such that for any given ordered pair of points there is a unique translation sending the first point to the second; the composition of two translations is their sum in the vector space of the translations. In more concrete terms, this amounts to having an operation that associates to any ordered pair of points a vector and another operation that allows translation of a point by a vector to give another point; these operations are required to satisfy a number of axioms (notably that two successive translations have the effect of translation by the sum vector). By choosing any point as "origin", the points are in one-to-one correspondence with the vectors, but there is no preferred choice for the origin; thus an affine space may be viewed as obtained from its associated vector space by "forgetting" the origin (zero vector). The idea of forgetting the metric can be applied in the theory of manifolds. That is developed in the article on the affine connection. History In 1748, Leonhard Euler introduced the term affine () in his book (volume 2, chapter XVIII). In 1827, August Möbius wrote on affine geometry in his (chapter 3). After Felix Klein's Erlangen program, affine geometry was recognized as a generalization of Euclidean geometry. In 1918, Hermann Weyl referred to affine geometry for his text Space, Time, Matter. He used affine geometry to introduce vector addition and subtraction at the earliest stages of his development of mathematical physics. Later, E. T. Whittaker wrote: Weyl's geometry is interesting historically as having been the first of the affine geometries to be worked out in detail: it is based on a special type of parallel transport [...using] worldlines of light-signals in four-dimensional space-time. A short element of one of these world-lines may be called a null-vector; then the parallel transport in question is such that it carries any null-vector at one point into the position of a null-vector at a neighboring point. Systems of axioms Several axiomatic approaches to affine geometry have been put forward: Pappus' law As affine geometry deals with parallel lines, one of the properties of parallels noted by Pappus of Alexandria has been taken as a premise: Suppose are on one line and on another. If the lines and are parallel and the lines and are parallel, then the lines and are parallel. (This is the affine version of Pappus's hexagon theorem). The full axiom system proposed has point, line, and line containing point as primitive notions: Two points are contained in just one line. For any line and any point , not on , there is just one line containing and not containing any point of . This line is said to be parallel to . Every line contains at least two points. There are at least three points not belonging to one line. According to H. S. M. Coxeter: The interest of these five axioms is enhanced by the fact that they can be developed into a vast body of propositions, holding not only in Euclidean geometry but also in Minkowski's geometry of time and space (in the simple case of 1 + 1 dimensions, whereas the special theory of relativity needs 1 + 3). The extension to either Euclidean or Minkowskian geometry is achieved by adding various further axioms of orthogonality, etc. The various types of affine geometry correspond to what interpretation is taken for rotation. Euclidean geometry corresponds to the ordinary idea of rotation, while Minkowski's geometry corresponds to hyperbolic rotation. With respect to perpendicular lines, they remain perpendicular when the plane is subjected to ordinary rotation. In the Minkowski geometry, lines that are hyperbolic-orthogonal remain in that relation when the plane is subjected to hyperbolic rotation. Ordered structure An axiomatic treatment of plane affine geometry can be built from the axioms of ordered geometry by the addition of two additional axioms: (Affine axiom of parallelism) Given a point and a line not through , there is at most one line through which does not meet . (Desargues) Given seven distinct points , such that are distinct lines through , and is parallel to , and is parallel to , then is parallel to . The affine concept of parallelism forms an equivalence relation on lines. Since the axioms of ordered geometry as presented here include properties that imply the structure of the real numbers, those properties carry over here so that this is an axiomatization of affine geometry over the field of real numbers. Ternary rings The first non-Desarguesian plane was noted by David Hilbert in his Foundations of Geometry. The Moulton plane is a standard illustration. In order to provide a context for such geometry as well as those where Desargues theorem is valid, the concept of a ternary ring was developed by Marshall Hall. In this approach affine planes are constructed from ordered pairs taken from a ternary ring. A plane is said to have the "minor affine Desargues property" when two triangles in parallel perspective, having two parallel sides, must also have the third sides parallel. If this property holds in the affine plane defined by a ternary ring, then there is an equivalence relation between "vectors" defined by pairs of points from the plane. Furthermore, the vectors form an abelian group under addition; the ternary ring is linear and satisfies right distributivity: Affine transformations Geometrically, affine transformations (affinities) preserve collinearity: so they transform parallel lines into parallel lines and preserve ratios of distances along parallel lines. We identify as affine theorems any geometric result that is invariant under the affine group (in Felix Klein's Erlangen programme this is its underlying group of symmetry transformations for affine geometry). Consider in a vector space , the general linear group . It is not the whole affine group because we must allow also translations by vectors in . (Such a translation maps any in to .) The affine group is generated by the general linear group and the translations and is in fact their semidirect product (Here we think of as a group under its operation of addition, and use the defining representation of on to define the semidirect product.) For example, the theorem from the plane geometry of triangles about the concurrence of the lines joining each vertex to the midpoint of the opposite side (at the centroid or barycenter) depends on the notions of mid-point and centroid as affine invariants. Other examples include the theorems of Ceva and Menelaus. Affine invariants can also assist calculations. For example, the lines that divide the area of a triangle into two equal halves form an envelope inside the triangle. The ratio of the area of the envelope to the area of the triangle is affine invariant, and so only needs to be calculated from a simple case such as a unit isosceles right angled triangle to give i.e. 0.019860... or less than 2%, for all triangles. Familiar formulas such as half the base times the height for the area of a triangle, or a third the base times the height for the volume of a pyramid, are likewise affine invariants. While the latter is less obvious than the former for the general case, it is easily seen for the one-sixth of the unit cube formed by a face (area 1) and the midpoint of the cube (height 1/2). Hence it holds for all pyramids, even slanting ones whose apex is not directly above the center of the base, and those with base a parallelogram instead of a square. The formula further generalizes to pyramids whose base can be dissected into parallelograms, including cones by allowing infinitely many parallelograms (with due attention to convergence). The same approach shows that a four-dimensional pyramid has 4D hypervolume one quarter the 3D volume of its parallelepiped base times the height, and so on for higher dimensions. Kinematics Two types of affine transformation are used in kinematics, both classical and modern. Velocity is described using length and direction, where length is presumed unbounded. This variety of kinematics, styled as Galilean or Newtonian, uses coordinates of absolute space and time. The shear mapping of a plane with an axis for each represents coordinate change for an observer moving with velocity in a resting frame of reference. Finite light speed, first noted by the delay in appearance of the moons of Jupiter, requires a modern kinematics. The method involves rapidity instead of velocity, and substitutes squeeze mapping for the shear mapping used earlier. This affine geometry was developed synthetically in 1912. to express the special theory of relativity. In 1984, "the affine plane associated to the Lorentzian vector space L2" was described by Graciela Birman and Katsumi Nomizu in an article entitled "Trigonometry in Lorentzian geometry". Affine space Affine geometry can be viewed as the geometry of an affine space of a given dimension n, coordinatized over a field K. There is also (in two dimensions) a combinatorial generalization of coordinatized affine space, as developed in synthetic finite geometry. In projective geometry, affine space means the complement of a hyperplane at infinity in a projective space. Affine space can also be viewed as a vector space whose operations are limited to those linear combinations whose coefficients sum to one, for example , , , , etc. Synthetically, affine planes are 2-dimensional affine geometries defined in terms of the relations between points and lines (or sometimes, in higher dimensions, hyperplanes). Defining affine (and projective) geometries as configurations of points and lines (or hyperplanes) instead of using coordinates, one gets examples with no coordinate fields. A major property is that all such examples have dimension 2. Finite examples in dimension 2 (finite affine planes) have been valuable in the study of configurations in infinite affine spaces, in group theory, and in combinatorics. Despite being less general than the configurational approach, the other approaches discussed have been very successful in illuminating the parts of geometry that are related to symmetry. Projective view In traditional geometry, affine geometry is considered to be a study between Euclidean geometry and projective geometry. On the one hand, affine geometry is Euclidean geometry with congruence left out; on the other hand, affine geometry may be obtained from projective geometry by the designation of a particular line or plane to represent the points at infinity. In affine geometry, there is no metric structure but the parallel postulate does hold. Affine geometry provides the basis for Euclidean structure when perpendicular lines are defined, or the basis for Minkowski geometry through the notion of hyperbolic orthogonality. In this viewpoint, an affine transformation is a projective transformation that does not permute finite points with points at infinity, and affine transformation geometry is the study of geometrical properties through the action of the group of affine transformations.	Mathematics	Non-Euclidean geometry	null
243916	https://en.wikipedia.org/wiki/Head%20louse	Head louse	The head louse (Pediculus humanus capitis) is an obligate ectoparasite of humans. Head lice are wingless insects that spend their entire lives on the human scalp and feed exclusively on human blood. Humans are the only known hosts of this specific parasite, while chimpanzees and bonobos host a closely related species, Pediculus schaeffi. Other species of lice infest most orders of mammals and all orders of birds. Lice differ from other hematophagic ectoparasites such as fleas in spending their entire lifecycle on a host. Head lice cannot fly, and their short, stumpy legs render them incapable of jumping, or even walking efficiently on flat surfaces. The non-disease-carrying head louse differs from the related disease-carrying body louse (Pediculus humanus humanus) in preferring to attach eggs to scalp hair rather than to clothing. The two subspecies are morphologically almost identical, but do not normally interbreed. From genetic studies, they are thought to have diverged as subspecies about 30,000–110,000 years ago, when many humans began to wear a significant amount of clothing. However, the degree of separation is contentious as they can produce fertile offspring in a laboratory. A much more distantly related species of hair-clinging louse, the pubic or crab louse (Pthirus pubis), also infests humans. It is morphologically different from the other two species and is much closer in appearance to the lice which infest other primates. Louse infestation of the body is known as pediculosis, pediculosis capitis for head lice, pediculosis corporis for body lice, and phthiriasis for pubic lice. Adult morphology Like other insects of the suborder Anoplura, adult head lice are small (2.5–3 mm long), dorsoventrally flattened (see anatomical terms of location), and wingless. The thoracic segments are fused, but otherwise distinct from the head and abdomen, the latter being composed of seven visible segments. Head lice are grey in general, but their precise color varies according to the environment in which they were raised. After feeding, consumed blood causes the louse body to take on a reddish color. Head One pair of antennae, each with five segments, protrudes from the insect's head. Head lice also have one pair of eyes. Eyes are present in all species within the Pediculidae family, but are reduced or absent in most other members of the Anoplura suborder. Like other members of the Anoplura, head louse mouthparts are highly adapted for piercing the skin and sucking blood. These mouth parts are retracted into the insect's head except during feeding. Thorax Six legs project from the fused segments of the thorax. As is typical in the Anoplura, these legs are short and terminate with a single claw and opposing "thumb". Between its claw and thumb, the louse grasps the hair of its host. With their short legs and large claws, lice are well adapted to clinging to the hair of their host. These adaptations leave them incapable of jumping, or even walking efficiently on flat surfaces. Lice can climb up strands of hair very quickly, allowing them to move quickly and reach another host. Abdomen Seven segments of the louse abdomen are visible. The first six segments each have a pair of spiracles through which the insect breathes. The last segment contains the anus and (separately) the genitalia. Sex differences In male lice, the front two legs are slightly larger than the other four. This specialized pair of legs are used for holding the female during copulation. Males are slightly smaller than females and are characterized by a pointed end of the abdomen and a well-developed genital apparatus visible inside the abdomen. Females are characterized by two gonopods in the shape of a W at the end of their abdomens. Eggs and nits Like most insects, head lice are oviparous. Females lay about three or four eggs per day. Louse eggs (also known as nits), are attached near the base of a host hair shaft. Eggs are usually laid on the base of the hair, 3–5 mm off the scalp surface. In warm climates, and especially the tropics, eggs may be laid or more down the hair shaft. To attach an egg, the adult female secretes a glue from her reproductive organ. This glue quickly hardens into a "nit sheath" that covers the hair shaft and large parts of the egg except for the operculum, a cap through which the embryo breathes. The glue was previously thought to be chitin-based, but more recent studies have shown it to be made of proteins similar to hair keratin. Each egg is oval-shaped and about 0.8 mm in length. They are bright, transparent, and tan to coffee-colored so long as they contain an embryo, but appear white after hatching. Head lice hatch typically six to nine days after oviposition. After hatching, the louse nymph leaves behind its egg shell, still attached to the hair shaft. The empty egg shell remains in place until physically removed by abrasion or the host, or until it slowly disintegrates, which may take six or more months. Empty shells are matte, collapsed and white in color. The term nit may include any of the following: Viable eggs that will eventually hatch Remnants of already-hatched eggs (nits) Nonviable eggs (dead embryo) that will never hatch Of these three, only eggs containing viable embryos have the potential to infest or reinfest a host. However, a no nit policy is a common public health measure to prevent transmission of lice. Some authors have therefore restricted the definition of nit to describe only a hatched or nonviable egg: Others have retained the broad definition, while simultaneously attempting to clarify its relevance to infestation: In British and Irish slang the term "nit" is often used, across different age groups, to refer to the head lice themselves. Development and nymphs Head lice, like other insects of the order Phthiraptera, are hemimetabolous. Newly hatched nymphs will moult three times before reaching the sexually mature adult stage. Thus, mobile head lice populations may contain eggs, nits, three nymphal instars, and the adults (male and female) (imago). Metamorphosis during head louse development is subtle. The only visible differences between different instars and the adult, other than size, is the relative length of the abdomen, which increases with each molt, as well as the existence of reproductive organs in the adults. Aside from reproduction, nymph behavior is similar to the adult. Like adults, nymphs feed also only on human blood (hematophagia), and cannot survive long away from a host. Outside their hosts lice cannot survive more than 24 hrs. The time required for head lice to complete their nymph development to the imago lasts for 12–15 days. Nymph mortality in captivity is about 38%, especially within the first two days of life. In the wild, mortality may instead be highest in the third instar. Nymph hazards are numerous. Failure to completely hatch from the egg is invariably fatal. Death during molting can also occur, although it is reportedly uncommon. During feeding, the nymph gut can rupture, dispersing the host's blood throughout the insect body. This results in death within a day or two. Whether the high mortality recorded under experimental conditions is representative of conditions in the wild is unclear. Reproduction and lifespan Head lice reproduce sexually, and copulation is necessary for the female to produce fertile eggs. Parthenogenesis, the production of viable offspring by virgin females, does not occur in Pediculus humanus. Pairing can begin within the first 10 hours of adult life. After 24 hours, adult lice copulate frequently, with mating occurring during any period of the night or day. Mating attachment frequently lasts more than an hour. Young males can successfully pair with older females, and vice versa. Experiments with P. h. humanus (body lice) emphasize the attendant hazards of lice copulation. A single young female confined with six or more males will die in a few days, having laid very few eggs. Similarly, death of a virgin female was reported after admitting a male to her confinement. The female laid only one egg after mating, and her entire body was tinged with red—a condition attributed to rupture of the alimentary canal during the sexual act. Old females frequently die following, if not during, copulation. During its lifespan of 4 weeks a female louse lays 50-150 eggs. Eggs hatch within 6–9 days, each nymphal stage last for 4–5 days and accordingly the period from egg to adults lasts for 18–24 days. Adult lice live for an additional 3–4 weeks. Factors affecting infestation The number of children per family, the sharing of beds and closets, hair washing habits, local customs and social contacts, healthcare in a particular area (e.g. school), and socioeconomic status were found to be significant factors in head louse infestation. Girls are two to four times more frequently infested than boys. Children between 4 and 14 years of age are the most frequently infested group. Behaviour Feeding All stages except eggs are blood-feeders and bite the skin four to five times daily to feed. They inject saliva which contains an anticoagulant and suck blood. The digested blood is excreted as dark red frass. Position on host Although any part of the scalp may be colonized, lice favor the nape of the neck and the area behind the ears, where the eggs are usually laid. Head lice are repelled by light and move towards shadows or dark-coloured objects in their vicinity. Transmission Lice have no wings or powerful legs for jumping, so they use the claws on their legs to move from hair to hair. Normally, head lice infest a new host only by close contact between individuals, making social contacts among children and parent-child interactions more likely routes of infestation than shared combs, hats, brushes, towels, clothing, beds, or closets. Head-to-head contact is by far the most common route of lice transmission. Distribution About 6–12 million people, mainly children, are treated annually for head lice in the United States alone. In the UK, it is estimated that two thirds of children will experience at least one case of head lice before leaving primary school. High levels of louse infestations have also been reported from all over the world, including Australia, Denmark, France, Ireland, Israel, and Sweden. Archaeogenetics Analysis of the DNA of lice found on Peruvian mummies may indicate that some diseases (such as typhus) may have passed from the New World to the Old World, instead of the other way around. Genome The sequencing of the genome of the body louse was first proposed in the mid-2000s and the annotated genome was published in 2010. An analysis of the body and head louse transcriptomes revealed these two organisms are extremely similar genetically. Unlike other bilateral animals, the 37 mitochondrial genes of human lice are not on a single circular chromosome but extensively fragmented. For the head louse, and the body louse, they are on 20 minichromosomes, for the pubic louse 14 minichromosomes and the chimpanzee louse,18 minichromosomes. Mitochondrial clades Human lice are divided into three deeply divergent mitochondrial clades known as A, B, and C. Three subclades have been identified, D (a sister clade of A), E (a sister clade of C), and F (a sister clade of B). Clade A head and body: worldwide found in ancient Roman Judea Clade D (sister of clade A) head and body: Central Africa, Ethiopia, United States Clade B head only: worldwide found in ancient Roman Judea and 4,000-year-old Chilean mummy Clade F (sister of clade B) head and body: South America Clade C head only: Ethiopia, Nepal, Thailand Clade E (sister of clade C) head only: West Africa	Biology and health sciences	Insects: General	Animals
243950	https://en.wikipedia.org/wiki/Godwit	Godwit	Godwits are a group of four large, long-billed, long-legged and strongly migratory waders of the bird genus Limosa. Their long bills allow them to probe deeply in the sand for aquatic worms and molluscs. In their winter range, they flock together where food is plentiful. They frequent tidal shorelines, breeding in northern climates in summer and migrating south in winter. A female bar-tailed godwit made a flight of 29,000 km (18,000 mi), flying of it without stopping. In 2020 a male bar-tailed godwit flew about non-stop in its migration from Alaska to New Zealand, previously a record for avian non-stop flight. In October 2022, a 5 month old, male bar-tailed godwit was tracked from Alaska to Tasmania, a trip that took 11 days, and recorded a non-stop flight of . The godwits can be distinguished from the curlews by their straight or slightly upturned bills, and from the dowitchers by their longer legs. The winter plumages are fairly drab, but three species have reddish underparts when breeding. The females are appreciably larger than the males. Godwits were once a popular British dish. Sir Thomas Browne writing in about 1682 noted that godwits "were accounted the daintiest dish in England". Taxonomy The genus Limosa was introduced by the French zoologist Mathurin Jacques Brisson in 1760 with the black-tailed godwit (Limosa limosa) as the type species. The genus name Limosa is from Latin and means "muddy", from limus, "mud". The English name "godwit" was first recorded in about 1416–17 and is believed to imitate the bird's call. The genus contains four living species: Fossil species In addition, there are two or three species of fossil prehistoric godwits. Limosa vanrossemi is known from the Monterey Formation (Late Miocene, approx. 6 mya) of Lompoc, United States. Limosa lacrimosa is known from the Early Pliocene of Western Mongolia (Kurochkin, 1985). Limosa gypsorum of the Late Eocene (Montmartre Formation, some 35 mya) of France may have actually been a curlew or some bird ancestral to both curlews and godwits (and possibly other Scolopacidae), or even a rail, being placed in the monotypic genus Montirallus by some (Olson, 1985). Certainly, curlews and godwits are rather ancient and in some respects primitive lineages of scolopacids, further complicating the assignment of such possibly basal forms. In a 2001 study comparing the ratios cerebrum to brain volumes in various dinosaur species, Hans C. E. Larsson found that more derived dinosaurs generally had proportionally more voluminous cerebrum. Limosa gypsorum, then regarded as a Numenius species, was a discrepancy in this general trend. L. gypsorum was only 63% of the way between a typical reptilian ratio and that of modern birds. However, this may be explainable if the endocast was distorted, as it had been previously depicted in the past by Deschaseaux, who is described by Larsson as calling the endocast "slightly anteroposteriorly sheared and laterally compressed." Citations General sources Larsson, H. C. E. 2001. Endocranial anatomy of Carcharodontosaurus saharicus (Theropoda: Allosauroidea) and its implications for theropod brain evolution. pp. 19–33. In: Mesozoic Vertebrate Life. Tanke, D. H., Carpenter, K., Skrepnick, M. W. (eds.). Indiana University Press. Olson, Storrs L. (1985): Section X.D.2.b. "Scolopacidae". In: Farner, D.S.; King, J.R. & Parkes, Kenneth C. (eds.): Avian Biology 8: 174–175. Academic Press, New York. Taxa named by Mathurin Jacques Brisson	Biology and health sciences	Charadriiformes	Animals
244091	https://en.wikipedia.org/wiki/Dowitcher	Dowitcher	The three dowitchers are medium-sized long-billed wading birds in the genus Limnodromus. The English name "dowitcher" is from Iroquois, recorded in English by the 1830s. They resemble godwits in body and bill shape, and the reddish underparts in summer, but are much shorter legged, more like snipes, to which they are more closely related. All three are strongly migratory. The two North American species are difficult to separate in most plumages, and were considered a single species for many years. The Asian bird is rare and not well known. Taxonomy The genus Limnodromus was introduced in 1833 by the German naturalist Prince Maximilian of Wied-Neuwied to accommodate a single species, the short-billed dowitcher. The name combines the Ancient Greek limnē meaning "marsh" with -dromos meaning "-racer" or "-runner". The dowitcher species are:	Biology and health sciences	Charadriiformes	Animals
244103	https://en.wikipedia.org/wiki/Tendinopathy	Tendinopathy	Tendinopathy is a type of tendon disorder that results in pain, swelling, and impaired function. The pain is typically worse with movement. It most commonly occurs around the shoulder (rotator cuff tendinitis, biceps tendinitis), elbow (tennis elbow, golfer's elbow), wrist, hip, knee (jumper's knee, popliteus tendinopathy), or ankle (Achilles tendinitis). Causes may include an injury or repetitive activities. Less common causes include infection, arthritis, gout, thyroid disease, diabetes and the use of quinolone antibiotic medicines. Groups at risk include people who do manual labor, musicians, and athletes. Diagnosis is typically based on symptoms, examination, and occasionally medical imaging. A few weeks following an injury little inflammation remains, with the underlying problem related to weak or disrupted tendon fibrils. Treatment may include rest, NSAIDs, splinting, and physiotherapy. Less commonly steroid injections or surgery may be done. About 80% of overuse tendinopathy patients recover completely within six months. Tendinopathy is relatively common. Older people are more commonly affected. It results in a large amount of missed work. Signs and symptoms Symptoms include tenderness on palpation, swelling, and pain, often when exercising or with a specific movement. Cause Causes may include an injury or repetitive activities. Groups at risk include people who do manual labor, musicians, and athletes. Less common causes include infection, arthritis, gout, thyroid disease, and diabetes. Successful treatments include rehabilitation therapy and/or surgery. Obesity, or more specifically, adiposity or fatness, has also been linked to an increasing incidence of tendinopathy. Quinolone antibiotics are associated with increased risk of tendinitis and tendon rupture. A 2013 review found the incidence of tendon injury among those taking fluoroquinolones to be between 0.08 and 0.2%. Fluoroquinolones most frequently affect large load-bearing tendons in the lower limb, especially the Achilles tendon which ruptures in approximately 30 to 40% of cases. Types Examples include: Achilles tendinitis Calcific tendinitis Patellar tendinitis (jumper's knee) Pathophysiology As of 2016, the pathophysiology of tendinopathy is poorly understood. While inflammation appears to play a role, the relationships among changes to the structure of tissue, the function of tendons, and pain are not understood and there are several competing models, none of which have been fully validated or falsified. Molecular mechanisms involved in inflammation includes release of inflammatory cytokines like IL-1β which reduces the expression of type I collagen mRNA in human tenocytes and causes extracellular matrix degradation in the tendon. In a 2020 systematic review, it was noted that while various inflammatory markers were present in two thirds of the reviewed articles, the heterogenicity of data and lack of comparable studies meant no conclusion about a common pathophysiology from this systematic review. There are multifactorial theories that could include: tensile overload, tenocyte related collagen synthesis disruption, load-induced ischemia, neural sprouting, thermal damage, and adaptive compressive responses. The intratendinous sliding motion of fascicles and shear force at interfaces of fascicles could be an important mechanical factor for the development of tendinopathy and predispose tendons to rupture. The most commonly accepted cause for this condition is seen to be an overuse syndrome in combination with intrinsic and extrinsic factors leading to what may be seen as a progressive interference or the failing of the innate healing response. Tendinopathy involves cellular apoptosis, matrix disorganization and neovascularization. Classic characteristics of "tendinosis" include degenerative changes in the collagenous matrix, hypercellularity, hypervascularity, and a lack of inflammatory cells which has challenged the original misnomer "tendinitis". For chronic tennis elbow, histological findings include granulation tissue, microrupture, degenerative changes, and there is no traditional inflammation. As a consequence, "lateral elbow tendinopathy or tendinosis" is used instead of "lateral epicondylitis". Examination of pathologic tennis elbow tissue reveals noninflammatory tissue, so the term "angiofibroblastic tendinosis" is also used. Cultures from tendinopathic tendons contain an increased production of type III collagen. Longitudinal sonogram of the lateral elbow displays thickening and heterogeneity of the common extensor tendon that is consistent with tendinosis, as the ultrasound reveals calcifications, intrasubstance tears, and marked irregularity of the lateral epicondyle. Although the term "epicondylitis" is frequently used to describe this disorder, most histopathologic findings of studies have displayed no evidence of an acute, or a chronic inflammatory process. Histologic studies have demonstrated that this condition is the result of tendon degeneration, which causes normal tissue to be replaced by a disorganized arrangement of collagen. Therefore, the disorder is more appropriately referred to as "tendinosis" or "tendinopathy" rather than "tendinitis". Colour Doppler ultrasound reveals structural tendon changes, with vascularity and hypo-echoic areas that correspond to the areas of pain in the extensor origin. Load-induced non-rupture tendinopathy in humans is associated with an increase in the ratio of collagen III:I proteins, a shift from large to small diameter collagen fibrils, buckling of the collagen fascicles in the tendon extracellular matrix, and buckling of the tenocyte cells and their nuclei. Diagnosis Symptoms can vary from aches or pains and local joint stiffness, to a burning that surrounds the whole joint around the inflamed tendon. In some cases, swelling occurs along with heat and redness, and there may be visible knots surrounding the joint. With this condition, the pain is usually worse during and after activity, and the tendon and joint area can become stiff the following day as muscles tighten from the movement of the tendon. Many patients report stressful situations in their life in correlation with the beginnings of pain which may contribute to the symptoms. Medical imaging Ultrasound imaging can be used to evaluate tissue strain, as well as other mechanical properties. Ultrasound-based techniques are becoming more popular because of its affordability, safety, and speed. Ultrasound can be used for imaging tissues, and the sound waves can also provide information about the mechanical state of the tissue. Treatment Treatment of tendon injuries is largely conservative. Use of non-steroidal anti-inflammatory drugs (NSAIDs), rest, and gradual return to exercise is a common therapy. A meta-analysis revealed that exercise using weights or a resistance band is more effective than using bodyweight alone. In addition, having rest days is more effective than exercising every day. Resting assists in the prevention of further damage to the tendon. Ice, compression and elevation are also frequently recommended. Physical therapy, occupational therapy, orthotics or braces may also be useful. Initial recovery is typically within two to three days and full recovery is within three to six months. Tendinosis occurs as the acute phase of healing has ended (six to eight weeks) but has left the area insufficiently healed. Treatment of tendinitis helps reduce some of the risks of developing tendinosis, which takes longer to heal. There is tentative evidence that low-level laser therapy may also be beneficial in treating tendinopathy. The effects of deep transverse friction massage for treating tennis elbow and lateral knee tendinitis is unclear. NSAIDs NSAIDs may be used to help with pain. They however do not alter long term outcomes. Other types of pain medication, like paracetamol (acetaminophen), may be just as useful. Steroids Steroid injections have not been shown to have long term benefits for tendonitis, but appear to improve pain and function in the short term more effectively than other treatments except NSAIDs. They appear to have little benefit in tendinitis of the rotator cuff. There are some concerns that they may have negative effects. Other injections There is insufficient evidence on the routine use of injection therapies (autologous blood, platelet-rich plasma, deproteinised haemodialysate, aprotinin, polysulphated glycosaminoglycan, skin derived fibroblasts etc.) for treating Achilles tendinopathy. As of 2014 there was insufficient evidence to support the use of platelet-rich therapies for treating musculoskeletal soft tissue injuries such as ligament, muscle and tendon tears and tendinopathies. Prognosis Initial recovery from overuse tendinosus is usually within two to three months, and 80% will recover fully within three to six months. Epidemiology Tendon injury and resulting tendinopathy are responsible for up to 30% of consultations to sports doctors and other musculoskeletal health providers. Tendinopathy is most often seen in tendons of athletes either before or after an injury but is becoming more common in non-athletes and sedentary populations. For example, the majority of patients with Achilles tendinopathy in a general population-based study did not associate their condition with a sporting activity. In another study the population incidence of Achilles tendinopathy increased sixfold from 1979–1986 to 1987–1994. The incidence of rotator cuff tendinopathy ranges from 0.3% to 5.5% and annual prevalence from 0.5% to 7.4%. Terminology Tendinitis is a very common, but misleading term. By definition, the suffix "-itis" means "inflammation of". Inflammation is the body's local response to tissue damage which involves red blood cells, white blood cells, blood proteins with dilation of blood vessels around the site of injury. Tendons are relatively avascular. Corticosteroids are drugs that reduce inflammation. Corticosteroids can be useful to relieve chronic tendinopathy pain, improve function, and reduce swelling in the short term. However, there is a greater risk of long-term recurrence. They are typically injected along with a small amount of a numbing drug called lidocaine. Research shows that tendons are weaker following corticosteroid injections. Tendinitis is still a very common diagnosis, though research increasingly documents that what is thought to be tendinitis is usually tendinosis. Anatomically close but separate conditions are: Enthesitis, wherein there is inflammation of the entheses, the sites where tendons or ligaments insert into the bone. It is associated with HLA B27 arthropathies such as ankylosing spondylitis, psoriatic arthritis, and reactive arthritis. Apophysitis, inflammation of the bony attachment, generally associated with overuse among growing children. Research The use of a nitric oxide delivery system (glyceryl trinitrate patches) applied over the area of maximal tenderness was found to reduce pain and increase range of motion and strength. A promising therapy involves eccentric loading exercises involving lengthening muscular contractions. Other animals Bowed tendon is a horseman's term for tendinitis (inflammation) and tendinosis (degeneration), most commonly seen in the superficial digital flexor tendon in the front leg of horses. Mesenchymal stem cells, derived from a horse's bone marrow or fat, are currently being used for tendon repair in horses.	Biology and health sciences	Types	Health
244113	https://en.wikipedia.org/wiki/Lyme%20disease	Lyme disease	Lyme disease, also known as Lyme borreliosis, is a tick-borne disease caused by species of Borrelia bacteria, transmitted by blood-feeding ticks in the genus Ixodes. The most common sign of infection is an expanding red rash, known as erythema migrans (EM), which appears at the site of the tick bite about a week afterwards. The rash is typically neither itchy nor painful. Approximately 70–80% of infected people develop a rash. Early diagnosis can be difficult. Other early symptoms may include fever, headaches and tiredness. If untreated, symptoms may include loss of the ability to move one or both sides of the face, joint pains, severe headaches with neck stiffness or heart palpitations. Months to years later, repeated episodes of joint pain and swelling may occur. Occasionally, shooting pains or tingling in the arms and legs may develop. Lyme disease is transmitted to humans by the bites of infected ticks of the genus Ixodes. In the United States, ticks of concern are usually of the Ixodes scapularis type. According to the Centers for Disease Control and Prevention, "In most cases, a tick must be attached for 36 to 48 hours or more before the Lyme disease bacterium can be transmitted. If you remove a tick quickly (within 24 hours), you can greatly reduce your chances of getting Lyme disease." In Europe, Ixodes ricinus ticks may spread the bacteria more quickly. In North America, the bacterial species Borrelia burgdorferi and B. mayonii cause Lyme disease. In Europe and Asia, Borrelia afzelii, Borrelia garinii, B. spielmanii and four other species also cause the disease. The disease does not appear to be transmissible between people, by other animals nor through food. Diagnosis is based on a combination of symptoms, history of tick exposure and possibly testing for specific antibodies in the blood. Blood tests are often falsely negative in the early stages of the disease. Testing of individual ticks is not typically useful. Prevention includes efforts to prevent tick bites by wearing clothing to cover the arms and legs and using DEET or picaridin-based insect repellents. Using pesticides to reduce tick numbers may also be effective. Ticks can be removed using tweezers. If the removed tick is full of blood a single dose of doxycycline may be used to prevent the development of infection but is not generally recommended since the development of infection is rare. If an infection develops, a number of antibiotics are effective, including doxycycline, amoxicillin and cefuroxime. Standard treatment usually lasts for two or three weeks. Some people develop a fever and muscle and joint pains from treatment, which may last for one or two days. In those who develop persistent symptoms, long-term antibiotic therapy has not been found to be useful. Lyme disease is the most common disease spread by ticks in the Northern Hemisphere. Infections are most common in the spring and early summer. Lyme disease was diagnosed as a separate condition for the first time in 1975 in Lyme, Connecticut. It was originally mistaken for juvenile rheumatoid arthritis. The bacterium involved was first described in 1981 by Willy Burgdorfer. People with persistent symptoms after appropriate treatments are said to have Post-Treatment Lyme Disease Syndrome (PTLDS). PTLDS is different from chronic Lyme disease, a term no longer supported by scientists and used in different ways by different groups. Some healthcare providers claim that chronic Lyme is caused by persistent infection, but this is not believed to be true because no evidence of persistent infection can be found after standard treatment. , clinical trials of proposed human vaccines for Lyme disease were being carried out, but no vaccine was available. A vaccine, LYMERix, was produced, but discontinued in 2002 due to insufficient demand. There are several vaccines for the prevention of Lyme disease in dogs. Signs and symptoms Lyme disease can produce a broad range of symptoms. The incubation period is usually one to two weeks, but can be much shorter (days) or much longer (months to years). Lyme symptoms most often occur from the month of May to September in the Northern Hemisphere because the nymphal stage of the tick is responsible for most cases. Early localized infection 80% of Lyme infections begin with a rash of some sort at the site of a tick bite, often near skin folds such as the armpit, groin, back of the knee, or the trunk under clothing straps, or in children's hair, ears, or neck. Most people who get infected do not remember seeing a tick or a bite. The rash appears typically one or two weeks (range 3–32 days) after the bite and expands 2–3 cm per day up to a diameter of 5–70 cm (median is 16 cm). The rash is usually circular or oval, red or bluish, and may have an elevated or darker center. This rash is termed an Erythema Migrans (EM) which translates as "Migrating Redness." In about 79% of cases in Europe, this rash gradually clears from the center toward the edges possibly forming a "bull's eye" or "target-like" pattern, but this clearing only happens in 19% of cases in endemic areas of the United States. The rash may feel warm, usually is not itchy, is rarely tender or painful, and takes up to four weeks to resolve if untreated. The Lyme rash is often accompanied by symptoms of a flu-like illness, including fatigue, headache, body aches, fever, and chills [though usually neither nausea nor upper-respiratory problems]. These symptoms may also appear without a rash or linger after the rash has disappeared. Lyme can progress to later stages without a rash or these symptoms. People with high fever for more than two days or whose other symptoms of viral-like illness do not improve despite antibiotic treatment for Lyme disease, or who have abnormally low levels of white or red cells or platelets in the blood, should be investigated for possible coinfection with other tick-borne diseases such as ehrlichiosis and babesiosis. Not everyone with Lyme disease has all the symptoms, and many of these symptoms can also occur with other diseases. Asymptomatic infection exists, but occurs in less than 7% of infected individuals in the United States. Asymptomatic infection may be much more common among those infected in Europe. Early disseminated infection Within days to weeks after the onset of local infection, the Borrelia bacteria may spread through the lymphatic system or bloodstream. In 10–20% of untreated cases, EM rashes develop at sites across the body that bear no relation to the original tick bite. Transient muscle pains and joint pains are also common. In about 10–15% of untreated people, Lyme causes neurological problems known as neuroborreliosis. Early neuroborreliosis typically appears 4–6 weeks (range 1–12 weeks) after the tick bite and involves some combination of lymphocytic meningitis, cranial neuritis, radiculopathy, and/or mononeuritis multiplex. Lymphocytic meningitis causes characteristic changes in the cerebrospinal fluid (CSF) and may be accompanied for several weeks by variable headache and, less commonly, usually mild meningitis signs such as inability to flex the neck fully and intolerance to bright lights but typically no or only very low fever. After several months neuroborreliosis can also present otolaryngological symptoms. Up to 76.5% of them present as tinnitus, the most common symptom. Vertigo and dizziness (53.7%) and hearing loss (16.7%) were the next most common symptoms. In children, partial loss of vision may also occur. Cranial neuritis is an inflammation of cranial nerves. When due to Lyme, it most typically causes facial palsy, impairing blinking, smiling, and chewing on one or both sides of the face. It may also cause intermittent double vision. Lyme radiculopathy is an inflammation of spinal nerve roots that often causes pain and less often weakness, numbness, or altered sensation in the areas of the body served by nerves connected to the affected roots, e.g. limb(s) or part(s) of trunk. The pain is often described as unlike any other previously felt, excruciating, migrating, worse at night, rarely symmetrical, and often accompanied by extreme sleep disturbance. Mononeuritis multiplex is an inflammation causing similar symptoms in one or more unrelated peripheral nerves. Rarely, early neuroborreliosis may involve inflammation of the brain or spinal cord, with symptoms such as confusion, abnormal gait, ocular movements, or speech, impaired movement, impaired motor planning, or shaking. In North America, facial palsy is the typical early neuroborreliosis presentation, occurring in 5–10% of untreated people, in about 75% of cases accompanied by lymphocytic meningitis. Lyme radiculopathy is reported half as frequently, but many cases may be unrecognized. In European adults, the most common presentation is a combination of lymphocytic meningitis and radiculopathy known as Bannwarth syndrome, accompanied in 36-89% of cases by facial palsy. In this syndrome, radicular pain tends to start in the same body region as the initial erythema migrans rash, if there was one, and precedes possible facial palsy and other impaired movement. In extreme cases, permanent impairment of motor or sensory function of the lower limbs may occur. In European children, the most common manifestations are facial palsy (in 55%), other cranial neuritis, and lymphocytic meningitis (in 27%). In about 4–10% of untreated cases in the United States and 0.3–4% of untreated cases in Europe, typically between June and December, about one month (range 4 days to 7 months) after the tick bite, the infection may cause heart complications known as Lyme carditis. Symptoms may include heart palpitations (in 69% of people), dizziness, fainting, shortness of breath, and chest pain. Other symptoms of Lyme disease may also be present, such as EM rash, joint aches, facial palsy, headaches, or radicular pain. In some people, however, carditis may be the first manifestation of Lyme disease. Lyme carditis in 19–87% of people adversely impacts the heart's electrical conduction system, causing atrioventricular block that often manifests as heart rhythms that alternate within minutes between abnormally slow and abnormally fast. In 10–15% of people, Lyme causes myocardial complications such as cardiomegaly, left ventricular dysfunction, or congestive heart failure. Another skin condition, found in Europe but not in North America, is borrelial lymphocytoma, a purplish lump that develops on the ear lobe, nipple, or scrotum. Late disseminated infection Lyme arthritis occurs in up to 60% of untreated people, typically starting about six months after infection. It usually affects only one or a few joints, often a knee or possibly the hip, other large joints, or the temporomandibular joint. Usually, large joint effusion and swelling occur, but only mild or moderate pain. Without treatment, swelling and pain typically resolve over time, but periodically return. Baker's cysts may form and rupture. In early US studies of Lyme disease, a rare peripheral neuropathy was described that included numbness, tingling, or burning starting at the feet or hands and over time possibly moving up the limbs. In a later analysis that discovered poor documentation of this manifestation, experts wondered if it exists at all in the US or is merely very rare. A neurologic syndrome called Lyme encephalopathy is associated with subtle memory and cognitive difficulties, insomnia, a general sense of feeling unwell, and changes in personality. Lyme encephalopathy is controversial in the US and has not been reported in Europe. Problems such as depression and fibromyalgia are as common in people with Lyme disease as in the general population. There is no compelling evidence that Lyme disease causes psychiatric disorders, behavioral disorders (e.g. ADHD), or developmental disorders (e.g. autism). Acrodermatitis chronica atrophicans is a chronic skin disorder observed primarily in Europe among the elderly. It begins as a reddish-blue patch of discolored skin, often on the backs of the hands or feet. The lesion slowly atrophies over several weeks or months, with the skin becoming first thin and wrinkled and then, if untreated, completely dry and hairless. It is also associated with peripheral neuropathy. Cause Lyme disease is caused by spirochetes, gram-negative bacteria from the genus Borrelia. Spirochetes are surrounded by peptidoglycan and flagella. The Lyme-related Borrelia species are collectively known as Borrelia burgdorferi sensu lato, and show a great deal of genetic diversity. B. burgdorferi sensu lato is a species complex made up of 20 accepted and three proposed genospecies. Eight species are known to cause Lyme disease: B. mayonii (found in North America), B. burgdorferi sensu stricto (found in North America and Europe), B. afzelii, B. garinii, B. spielmanii, and B. lusitaniae (all found in Eurasia). Some studies have also proposed that B. valaisiana may sometimes infect humans, but this species does not seem to be an important cause of disease. Tick life cycle Three stages occur in the life cycle of a tick - larva, nymph, and adult. During the nymph stage, ticks most frequently transmit Lyme disease and are usually most active in late spring and early summer in regions where the climate is mild. During the adult stage, Lyme disease transmission is less common because adult ticks are less likely to bite humans and tend to be larger in size, so can be easily seen and removed. Transmission Lyme disease is classified as a zoonosis, as it is transmitted to humans from a natural reservoir among small mammals and birds by ticks that feed on both sets of hosts. Hard-bodied ticks of the genus Ixodes are the vectors of Lyme disease (also the vector for Babesia). Most infections are caused by ticks in the nymphal stage, because they are very small, thus may feed for long periods of time undetected. Nymphal ticks are generally the size of a poppy seed and sometimes with a dark head and a translucent body. Or, the nymphal ticks can be darker. The younger larval ticks are very rarely infected. Although deer are the preferred hosts of adult deer ticks, and tick populations are much lower in the absence of deer, ticks generally do not acquire Borrelia from deer, instead they obtain them from infected small mammals such as the white-footed mouse, and occasionally birds. Areas where Lyme is common are expanding. Within the tick midgut, the Borrelias outer surface protein A (OspA) binds to the tick receptor for OspA, known as TROSPA. When the tick feeds, the Borrelia downregulates OspA and upregulates OspC, another surface protein. After the bacteria migrate from the midgut to the salivary glands, OspC binds to Salp15, a tick salivary protein that appears to have immunosuppressive effects that enhance infection. Successful infection of the mammalian host depends on bacterial expression of OspC. Tick bites often go unnoticed because of the small size of the tick in its nymphal stage, as well as tick secretions that prevent the host from feeling any itch or pain from the bite. However, transmission is quite rare, with only about 1.2 to 1.4 percent of recognized tick bites resulting in Lyme disease. While B. burgdorferi is most associated with ticks hosted by white-tailed deer and white-footed mice, Borrelia afzelii is most frequently detected in rodent-feeding vector ticks, and Borrelia garinii and Borrelia valaisiana appear to be associated with birds. Both rodents and birds are competent reservoir hosts for B. burgdorferi sensu stricto. The resistance of a genospecies of Lyme disease spirochetes to the bacteriolytic activities of the alternative complement pathway of various host species may determine its reservoir host association. Budding research has suggested that B. burgdorferi sensu lato may also be able to form enzootic cycle among lizard populations; this was previously assumed not to be possible in major areas containing populations of lizards, such as California. Except for one study in Europe, much of the data implicating lizards is based on DNA detection of the spirochete and has not demonstrated that lizards are able to infect ticks feeding upon them. As some experiments suggest lizards are refractory to infection with Borrelia, it appears likely their involvement in the enzootic cycle is more complex and species-specific. In Europe, the main vector is Ixodes ricinus, which is also called the sheep tick or castor bean tick. In China, Ixodes persulcatus (the taiga tick) is probably the most important vector. In North America, the black-legged tick or deer tick (Ixodes scapularis) is the main vector on the East Coast. The lone star tick (Amblyomma americanum), which is found throughout the Southeastern United States as far west as Texas, is unlikely to transmit the Lyme disease spirochetes, though it may be implicated in a related syndrome called southern tick-associated rash illness, which resembles a mild form of Lyme disease. On the West Coast of the United States, the main vector is the western black-legged tick (Ixodes pacificus). The tendency of this tick species to feed predominantly on host species such as the Western Fence Lizard that are resistant to Borrelia infection appears to diminish transmission of Lyme disease in the West. Transmission can occur across the placenta during pregnancy and as with a number of other spirochetal diseases, adverse pregnancy outcomes are possible with untreated infection; prompt treatment with antibiotics reduces or eliminates this risk. There is no scientific evidence to support Lyme disease transmission via blood transfusion, sexual contact, or breast milk. Tick-borne co-infections Ticks that transmit B. burgdorferi to humans can also carry and transmit several other microbes, such as Babesia microti and Anaplasma phagocytophilum, which cause the diseases babesiosis and human granulocytic anaplasmosis (HGA), respectively. Among people with early Lyme disease, depending on their location, 2–12% will also have HGA and 2–10% will have babesiosis. Ticks in certain regions also transmit viruses that cause tick-borne encephalitis and Powassan virus disease. Co-infections of Lyme disease may not require additional treatment, since they may resolve on their own or—as in the case of HGA—can be treated with the doxycycline prescribed for Lyme. Persistent fever or compatible anomalous laboratory findings may be indicative of a co-infection. Pathophysiology B. burgdorferi can spread throughout the body during the course of the disease, and has been found in the skin, heart, joints, peripheral nervous system, and central nervous system. B. Burgdorferi does not produce toxins. Therefore, many of the signs and symptoms of Lyme disease are a consequence of the immune response to spirochete in those tissues. B. burgdorferi is injected into the skin by the bite of an infected Ixodes tick. Tick saliva, which accompanies the spirochete into the skin during the feeding process, contains substances that disrupt the immune response at the site of the bite. This provides a protective environment where the spirochete can establish infection. The spirochetes multiply and migrate outward within the dermis. The host inflammatory response to the bacteria in the skin causes the characteristic circular EM lesion. Neutrophils, however, which are necessary to eliminate the spirochetes from the skin, fail to appear in necessary numbers in the developing EM lesion because tick saliva inhibits neutrophil function. This allows the bacteria to survive and eventually spread throughout the body. Days to weeks following the tick bite, the spirochetes spread via the bloodstream to joints, heart, nervous system, and distant skin sites, where their presence gives rise to the variety of symptoms of the disseminated disease. The spread of B. burgdorferi is aided by the attachment of the host protease plasmin to the surface of the spirochete. If untreated, the bacteria may persist in the body for months or even years, despite the production of B. burgdorferi antibodies by the immune system. The spirochetes may avoid the immune response by decreasing expression of surface proteins that are targeted by antibodies, antigenic variation of the VlsE surface protein, inactivating key immune components such as complement, and hiding in the extracellular matrix, which may interfere with the function of immune factors. Immunological studies Exposure to the Borrelia bacterium during Lyme disease possibly causes a long-lived and damaging inflammatory response, a form of pathogen-induced autoimmune disease. The production of this reaction might be due to a form of molecular mimicry, where Borrelia avoids being killed by the immune system by resembling normal parts of the body's tissues. Chronic symptoms from an autoimmune reaction could explain why some symptoms persist even after the spirochetes have been eliminated from the body. This hypothesis may explain why chronic arthritis persists after antibiotic therapy, similar to rheumatic fever, but its wider application is controversial. Diagnosis Lyme disease is diagnosed based on symptoms, objective physical findings (such as erythema migrans (EM) rash, facial palsy, or arthritis), history of possible exposure to infected ticks, and possibly laboratory tests. People with symptoms of early Lyme disease should have a total body skin examination for EM rashes and asked whether EM-type rashes had manifested within the last 1–2 months. Presence of an EM rash and recent tick exposure (i.e., being outdoors in a likely tick habitat where Lyme is common, within 30 days of the appearance of the rash) are sufficient for Lyme diagnosis; no laboratory confirmation is needed or recommended. Most people who get infected do not remember a tick or a bite, and the EM rash need not look like a bull's eye (most EM rashes in the U.S. do not) or be accompanied by any other symptoms. In the U.S., Lyme is most common in the New England and Mid-Atlantic states and parts of Wisconsin and Minnesota, but it is expanding into other areas. Several bordering areas of Canada also have high Lyme risk. In the absence of an EM rash or history of tick exposure, Lyme diagnosis depends on laboratory confirmation. The bacteria that cause Lyme disease are difficult to observe directly in body tissues and also difficult and too time-consuming to grow in the laboratory. The most widely used tests look instead for presence of antibodies against those bacteria in the blood. A positive antibody test result does not by itself prove active infection but can confirm an infection that is suspected because of symptoms, objective findings, and history of tick exposure in a person. Because as many as 5–20% of the normal population have antibodies against Lyme, people without history and symptoms suggestive of Lyme disease should not be tested for Lyme antibodies: a positive result would likely be false, possibly causing unnecessary treatment. In some cases, when history, signs, and symptoms are strongly suggestive of early disseminated Lyme disease, empiric treatment may be started and reevaluated as laboratory test results become available. Laboratory testing Tests for antibodies in the blood by ELISA and Western blot is the most widely used method for Lyme diagnosis. A two-tiered protocol is recommended by the Centers for Disease Control and Prevention (CDC): the sensitive ELISA test is performed first, and if it is positive or equivocal, then the more specific Western blot is run. The immune system takes some time to produce antibodies in quantity. After Lyme infection onset, antibodies of types IgM and IgG usually can first be detected respectively at 2–4 weeks and 4–6 weeks, and peak at 6–8 weeks. When an EM rash first appears, detectable antibodies may not be present. Therefore, it is recommended that testing not be performed and diagnosis be based on the presence of the EM rash. Up to 30 days after suspected Lyme infection onset, infection can be confirmed by detection of IgM or IgG antibodies; after that, it is recommended that only IgG antibodies be considered. A positive IgM and negative IgG test result after the first month of infection is generally indicative of a false-positive result. The number of IgM antibodies usually collapses 4–6 months after infection, while IgG antibodies can remain detectable for years. Other tests may be used in neuroborreliosis cases. In Europe, neuroborreliosis is usually caused by Borrelia garinii and almost always involves lymphocytic pleocytosis, i.e. the densities of lymphocytes (infection-fighting cells) and protein in the cerebrospinal fluid (CSF) typically rise to characteristically abnormal levels, while glucose level remains normal. Additionally, the immune system produces antibodies against Lyme inside the intrathecal space, which contains the CSF. Demonstration by lumbar puncture and CSF analysis of pleocytosis and intrathecal antibody production are required for definite diagnosis of neuroborreliosis in Europe (except in cases of peripheral neuropathy associated with acrodermatitis chronica atrophicans, which usually is caused by Borrelia afzelii and confirmed by blood antibody tests). In North America, neuroborreliosis is caused by Borrelia burgdorferi and may not be accompanied by the same CSF signs; they confirm a diagnosis of central nervous system (CNS) neuroborreliosis if positive, but do not exclude it if negative. American guidelines consider CSF analysis optional when symptoms appear to be confined to the peripheral nervous system (PNS), e.g. facial palsy without overt meningitis symptoms. Unlike blood and intrathecal antibody tests, CSF pleocytosis tests revert to normal after infection ends and therefore can be used as objective markers of treatment success and inform decisions on whether to retreat. In infection involving the PNS, electromyography and nerve conduction studies can be used to monitor objectively the response to treatment. In Lyme carditis, electrocardiograms are used to evidence heart conduction abnormalities, while echocardiography may show myocardial dysfunction. Biopsy and confirmation of Borrelia cells in myocardial tissue may be used in specific cases but are usually not done because of risk of the procedure. Polymerase chain reaction (PCR) tests for Lyme disease have also been developed to detect the genetic material (DNA) of the Lyme disease spirochete. Culture or PCR are the current means for detecting the presence of the organism, as serologic studies only test for antibodies of Borrelia. PCR has the advantage of being much faster than culture. However, PCR tests are susceptible to false positive results, e.g. by detection of debris of dead Borrelia cells or specimen contamination. Even when properly performed, PCR often shows false-negative results because few Borrelia cells can be found in blood and cerebrospinal fluid (CSF) during infection. Hence, PCR tests are recommended only in special cases, e.g. diagnosis of Lyme arthritis, because it is a highly sensitive way of detecting ospA DNA in synovial fluid. Although sensitivity of PCR in CSF is low, its use may be considered when intrathecal antibody production test results are suspected of being falsely negative, e.g. in very early (< 6 weeks) neuroborreliosis or in immunosuppressed people. Several other forms of laboratory testing for Lyme disease are available, some of which have not been adequately validated. OspA antigens, shed by live Borrelia bacteria into urine, are a promising technique being studied. The use of nanotrap particles for their detection is being looked at and the OspA has been linked to active symptoms of Lyme. High titers of either immunoglobulin G (IgG) or immunoglobulin M (IgM) antibodies to Borrelia antigens indicate disease, but lower titers can be misleading, because the IgM antibodies may remain after the initial infection, and IgG antibodies may remain for years. The CDC does not recommend urine antigen tests, PCR tests on urine, immunofluorescent staining for cell-wall-deficient forms of B. burgdorferi, and lymphocyte transformation tests. Imaging Neuroimaging is controversial in whether it provides specific patterns unique to neuroborreliosis, but may aid in differential diagnosis and in understanding the pathophysiology of the disease. Though controversial, some evidence shows certain neuroimaging tests can provide data that are helpful in the diagnosis of a person. Magnetic resonance imaging (MRI) and single-photon emission computed tomography (SPECT) are two of the tests that can identify abnormalities in the brain of a person affected with this disease. Neuroimaging findings in an MRI include lesions in the periventricular white matter, as well as enlarged ventricles and cortical atrophy. The findings are considered somewhat unexceptional because the lesions have been found to be reversible following antibiotic treatment. Images produced using SPECT show numerous areas where an insufficient amount of blood is being delivered to the cortex and subcortical white matter. However, SPECT images are known to be nonspecific because they show a heterogeneous pattern in the imaging. The abnormalities seen in the SPECT images are very similar to those seen in people with cerebral vasculitis and Creutzfeldt–Jakob disease, which makes them questionable. Differential diagnosis Community clinics have been reported to misdiagnose 23–28% of Erythema migrans (EM) rashes and 83% of other objective manifestations of early Lyme disease. EM rashes are often misdiagnosed as spider bites, cellulitis, or shingles. Many misdiagnoses are credited to the widespread misconception that EM rashes should look like a bull's eye. Actually, the key distinguishing features of the EM rash are the speed and extent to which it expands, respectively up to 2–3 cm/day and a diameter of at least 5 cm, and in 50% of cases more than 16 cm. The rash expands away from its center, which may or may not look different or be separated by ring-like clearing from the rest of the rash. Compared to EM rashes, spider bites are more common in the limbs, tend to be more painful and itchy or become swollen, and some may cause necrosis (sinking dark blue patch of dead skin). Cellulitis most commonly develops around a wound or ulcer, is rarely circular, and is more likely to become swollen and tender. EM rashes often appear at sites that are unusual for cellulitis, such as the armpit, groin, abdomen, or back of knee. Like Lyme, shingles often begins with headache, fever, and fatigue, which are followed by pain or numbness. However, unlike Lyme, in shingles these symptoms are usually followed by appearance of rashes composed of multiple small blisters along with a nerve's dermatome, and shingles can also be confirmed by quick laboratory tests. Facial palsy caused by Lyme disease (LDFP) is often misdiagnosed as Bell's palsy. Although Bell's palsy is the most common type of one-sided facial palsy (about 70% of cases), LDFP can account for about 25% of cases of facial palsy in areas where Lyme disease is common. Compared to LDFP, Bell's palsy much less frequently affects both sides of the face. Even though LDFP and Bell's palsy have similar symptoms and evolve similarly if untreated, corticosteroid treatment is beneficial for Bell's Palsy, while being detrimental for LDFP. Recent history of exposure to a likely tick habitat during warmer months, EM rash, viral-like symptoms such as headache and fever, and/or palsy in both sides of the face should be evaluated for the likelihood of LDFP; if it is more than minimal, empiric therapy with antibiotics should be initiated, without corticosteroids, and reevaluated upon completion of laboratory tests for Lyme disease. Unlike viral meningitis, Lyme lymphocytic meningitis tends to not cause fever, last longer, and recur. Lymphocytic meningitis is also characterized by possibly co-occurring with EM rash, facial palsy, or partial vision obstruction and having much lower percentage of polymorphonuclear leukocytes in CSF. Lyme radiculopathy affecting the limbs is often misdiagnosed as a radiculopathy caused by nerve root compression, such as sciatica. Although most cases of radiculopathy are compressive and resolve with conservative treatment (e.g., rest) within 4–6 weeks, guidelines for managing radiculopathy recommend first evaluating risks of other possible causes that, although less frequent, require immediate diagnosis and treatment, including infections such as Lyme and shingles. A history of outdoor activities in likely tick habitats in the last 3 months possibly followed by a rash or viral-like symptoms, and current headache, other symptoms of lymphocytic meningitis, or facial palsy would lead to suspicion of Lyme disease and recommendation of serological and lumbar puncture tests for confirmation. Lyme radiculopathy affecting the trunk can be misdiagnosed as myriad other conditions, such as diverticulitis and acute coronary syndrome. Diagnosis of late-stage Lyme disease is often complicated by a multifaceted appearance and nonspecific symptoms, prompting one reviewer to call Lyme the new "great imitator". As all people with later-stage infection will have a positive antibody test, simple blood tests can exclude Lyme disease as a possible cause of a person's symptoms. Prevention Tick bites may be prevented by avoiding or reducing time in likely tick habitats and taking precautions while in and when getting out of one. Most Lyme human infections are caused by Ixodes nymph bites between April and September. Ticks prefer moist, shaded locations in woodlands, shrubs, tall grasses and leaf litter or wood piles. Tick densities tend to be highest in woodlands, followed by unmaintained edges between woods and lawns (about half as high), ornamental plants and perennial groundcover (about a quarter), and lawns (about 30 times less). Ixodes larvae and nymphs tend to be abundant also where mice nest, such as stone walls and wood logs. Ixodes larvae and nymphs typically wait for potential hosts ("quest") on leaves or grasses close to the ground with forelegs outstretched; when a host brushes against its limbs, the tick rapidly clings and climbs on the host looking for a skin location to bite. In Northeastern United States, 69% of tick bites are estimated to happen in residences, 11% in schools or camps, 9% in parks or recreational areas, 4% at work, 3% while hunting, and 4% in other areas. Activities associated with tick bites around residences include yard work, brush clearing, gardening, playing in the yard, and letting dogs or cats that roam outside in woody or grassy areas into the house. In parks, tick bites often happen while hiking or camping. Walking on a mown lawn or center of a trail without touching adjacent vegetation is less risky than crawling or sitting on a log or stone wall. Pets should not be allowed to roam freely in likely tick habitats. As a precaution, CDC recommends soaking or spraying clothes, shoes, and camping gear such as tents, backpacks and sleeping bags with 0.5% permethrin solution and hanging them to dry before use. Permethrin is odorless and safe for humans but highly toxic to ticks. After crawling on permethrin-treated fabric for as few as 10–20 seconds, tick nymphs become irritated and fall off or die. Permethrin-treated closed-toed shoes and socks reduce by 74 times the number of bites from nymphs that make first contact with a shoe of a person also wearing treated shorts (because nymphs usually quest near the ground, this is a typical contact scenario). Better protection can be achieved by tucking permethrin-treated trousers (pants) into treated socks and a treated long-sleeve shirt into the trousers so as to minimize gaps through which a tick might reach the wearer's skin. Light-colored clothing may make it easier to see ticks and remove them before they bite. Military and outdoor workers' uniforms treated with permethrin have been found to reduce the number of bite cases by 80–95%. Permethrin protection lasts several weeks of wear and washings in customer-treated items and up to 70 washings for factory-treated items. Permethrin should not be used on human skin, underwear or cats. The EPA recommends several tick repellents for use on exposed skin, including DEET, picaridin, IR3535 (a derivative of amino acid beta-alanine), oil of lemon eucalyptus (OLE, a natural compound) and OLE's active ingredient para-menthane-diol (PMD). Unlike permethrin, repellents repel but do not kill ticks, protect for only several hours after application, and may be washed off by sweat or water. The most popular repellent is DEET in the U.S. and picaridin in Europe. Unlike DEET, picaridin is odorless and is less likely to irritate the skin or harm fabric or plastics. Repellents with higher concentration may last longer but are not more effective; against ticks, 20% picaridin may work for 8 hours vs. 55–98.11% DEET for 5–6 hours or 30–40% OLE for 6 hours. Repellents should not be used under clothes, on eyes, mouth, wounds or cuts, or on babies younger than 2 months (3 years for OLE or PMD). If sunscreen is used, repellent should be applied on top of it. Repellents should not be sprayed directly on a face, but should instead be sprayed on a hand and then rubbed on the face. After coming indoors, clothes, gear and pets should be checked for ticks. Clothes can be put into a hot dryer for 10 minutes to kill ticks (just washing or warm dryer are not enough). Showering as soon as possible, looking for ticks over the entire body, and removing them reduce risk of infection. Unfed tick nymphs are the size of a poppy seed, but a day or two after biting and attaching themselves to a person, they look like a small blood blister. The following areas should be checked especially carefully: armpits, between legs, back of knee, bellybutton, trunk, and in children ears, neck and hair. Tick removal Attached ticks should be removed promptly. Risk of infection increases with time of attachment, but in North America risk of Lyme disease is small if the tick is removed within 36 hours. CDC recommends inserting a fine-tipped tweezer between the skin and the tick, grasping very firmly, and pulling the closed tweezer straight away from the skin without twisting, jerking, squeezing or crushing the tick. After tick removal, any tick parts remaining in the skin should be removed with a clean tweezer, if possible. The wound and hands should then be cleaned with alcohol or soap and water. The tick may be disposed by placing it in a container with alcohol, sealed bag, tape or flushed down the toilet. The bitten person should write down where and when the bite happened so that this can be informed to a doctor if the person gets a rash or flu-like symptoms in the following several weeks. CDC recommends not using fingers, nail polish, petroleum jelly or heat on the tick to try to remove it. In Australia, where the Australian paralysis tick is prevalent, the Australasian Society of Clinical Immunology and Allergy recommends not using tweezers to remove ticks, because if the person is allergic, anaphylaxis could result. Instead, a product should be sprayed on the tick to cause it to freeze and then drop off. Another method consists in using about 20 cm of dental floss or fishing line for slowly tying an overhand knot between the skin and the tick and then pulling it away from the skin. Preventive antibiotics The risk of infectious transmission increases with the duration of tick attachment. It requires between 36 and 48 hours of attachment for the bacteria that causes Lyme to travel from within the tick into its saliva. If a deer tick that is sufficiently likely to be carrying Borrelia is found attached to a person and removed, and if the tick has been attached for 36 hours or is engorged, a single dose of doxycycline administered within the 72 hours after removal may reduce the risk of Lyme disease. It is not generally recommended for all people bitten, as development of infection is rare: about 50 bitten people would have to be treated this way to prevent one case of erythema migrans (i.e. the typical rash found in about 70–80% of people infected). Garden landscaping Several landscaping practices may reduce the risk of tick bites in residential yards. These include keeping lawns mown, removing leaf litter and weeds and avoiding the use of ground cover. A 3-ft-wide rock or woodchip barrier is recommended to separate lawns from wood piles, woodlands, stone walls and shrubs. Without vegetation on the barrier, ticks will tend not to cross it; acaricides may also be sprayed on it to kill ticks. A sun-exposed tick-safe zone at least 9 ft from the barrier should concentrate human activity on the yard, including any patios, playgrounds and gardening. Materials such as wood decking, concrete, bricks, gravel or woodchips used on the ground under patios and playgrounds would discourage ticks there. An 8-ft-high fence may be added to keep deer away from the tick-safe zone. Occupational exposure Outdoor workers are at risk of Lyme disease if they work at sites with infected ticks. This includes construction, landscaping, forestry, brush clearing, land surveying, farming, railroad work, oil field work, utility line work, park or wildlife management. U.S. workers in the northeastern and north-central states are at highest risk of exposure to infected ticks. Ticks may also transmit other tick-borne diseases to workers in these and other regions of the country. Worksites with woods, bushes, high grass or leaf litter are likely to have more ticks. Outdoor workers should be most careful to protect themselves in the late spring and summer when young ticks are most active. Host animals Ticks can feed upon the blood of a wide array of possible host species, including lizards, birds, mice, cats, dogs, deer, cattle and humans. The extent to which a tick can feed, reproduce, and spread will depend on the type and availability of its hosts. Whether it will spread disease is also affected by its available hosts. Some species, such as lizards, are referred to as "dilution hosts" because they don't tend to support Lyme disease pathogens and so decrease the likelihood that the disease will be passed on by ticks feeding on them. White-tailed deer are both a food source and a "reproductive host", where ticks tend to mate. The white-footed mouse is a reservoir host in which the pathogen for Lyme disease can survive. Availability of hosts can have significant impacts on the transmission of Lyme disease. A greater diversity of hosts, or of those that don't support the pathogen, tends to decrease the likelihood that the disease will be transmitted. In the United States, one approach to reducing the incidence of Lyme and other deer tick-borne diseases has been to greatly reduce the deer population on which the adult ticks depend for feeding and reproduction. Lyme disease cases fell following deer eradication on an island, Monhegan, Maine, and following deer control in Mumford Cove, Connecticut. Advocates have suggested reducing the deer population to levels of 8 to 10 deer per square mile, compared to levels of 60 or more deer per square mile in the areas of the country with the highest Lyme disease rates. While these studies have found these effects, other studies have found opposite effects. A study done in Massachusetts removed deer and did not see a significant decrease in tick abundance afterwards. Another study done in New Jersey removed deer and also did not see a reduction in the number of questing ticks and determined that deer culling is an unlikely way to effectively control tick populations. One study summarized the results of multiple studies all looking at deer reduction controlling tick populations and determined that deer control can't be used as a standalone reduction for Lyme disease. It also claims that since most of the studies looking at this are not good representatives of areas with high human Lyme disease risk. There is varying information on whether or not the removal of deer is actually a way to control the Lyme disease epidemic. Removal of smaller mammals that are fed on by juveniles who are more actively acquiring and spreading the pathogen, would decrease Lyme disease risk the most. Others have noted that while deer are reproductive hosts, they are not Borrelia burgdorferi reservoirs. This is because it was found that white-tailed deer blood actually kills the Borrelia burgdorferi bacteria. Researchers have suggested that smaller, less obviously visible Lyme reservoirs, like white-footed mice and Eastern chipmunks, may more strongly impact Lyme disease occurrence. Ecosystem studies in New York state suggest that white-footed mice thrive when forests are broken into smaller isolated chunks of woodland with fewer rodent predators. With more rodents harboring the disease, the odds increase that a tick will feed on a disease-harboring rodent and that someone will pick up a disease-carrying tick in their garden or walking in the woods. Data indicates that the smaller the wooded area, the more ticks it will contain and the likely they are to carry Lyme disease, supporting the idea that deforestation and habitat fragmentation affect ticks, hosts and disease transmission. Tick-borne diseases are estimated to affect ~80 % of cattle worldwide. They also affect cats, dogs, and other pets. Routine veterinary control of ticks of domestic animals through the use of acaricides has been suggested as a way to reduce exposure of humans to ticks. However, chemical control with acaricides is now criticized on a number of grounds. Ticks appear to develop resistance to acaricides; acaricides are costly; and there are concerns over their toxicity and the potential for chemical residues to affect food and the environment. In Europe, known reservoirs of Borrelia burgdorferi were 9 small mammals, 7 medium-sized mammals and 16 species of birds (including passerines, sea-birds and pheasants). These animals seem to transmit spirochetes to ticks and thus participate in the natural circulation of B. burgdorferi in Europe. The house mouse is also suspected as well as other species of small rodents, particularly in Eastern Europe and Russia. "The reservoir species that contain the most pathogens are the European roe deer Capreolus capreolus; "it does not appear to serve as a major reservoir of B. burgdorferi" thought Jaenson & al. (1992) (incompetent host for B. burgdorferi and TBE virus) but it is important for feeding the ticks, as red deer and wild boars (Sus scrofa), in which one Rickettsia and three Borrelia species were identified", with high risks of coinfection in roe deer. Nevertheless, in the 2000s, in roe deer in Europe "two species of Rickettsia and two species of Borrelia were identified". Vaccination no human vaccines for Lyme disease were available. The only human vaccine to advance to market was LYMErix, which was available from 1998, but discontinued in 2002. The vaccine candidate VLA15 was scheduled to start a phase 3 trial in the third quarter of 2022, with other research ongoing. Multiple vaccines are available for the prevention of Lyme disease in dogs. LYMErix The vaccine LYMErix was available from 1998 to 2002. The recombinant vaccine against Lyme disease, based on the outer surface protein A (OspA) of B. burgdorferi with aluminum hydroxide as adjuvant, was developed by SmithKline Beecham. In clinical trials involving more than 10,000 people, the vaccine was found to confer protective immunity to Lyme disease in 76% of adults after three doses with only mild or moderate and transient adverse effects. On 21 December 1998, the Food and Drug Administration (FDA) approved LYMErix on the basis of these trials for persons of ages 15 through 70. Following approval of the vaccine, its entry into clinical practice was slow for a variety of reasons, including its cost, which was often not reimbursed by insurance companies. Subsequently, hundreds of vaccine recipients reported they had developed autoimmune and other side effects. Supported by some advocacy groups, a number of class-action lawsuits were filed against GlaxoSmithKline, alleging the vaccine had caused these health problems. These claims were investigated by the FDA and the Centers for Disease Control, which found no connection between the vaccine and the autoimmune complaints. Despite the lack of evidence that the complaints were caused by the vaccine, sales plummeted and LYMErix was withdrawn from the U.S. market by GlaxoSmithKline in February 2002, in the setting of negative media coverage and fears of vaccine side effects. The fate of LYMErix was described in the medical literature as a "cautionary tale"; an editorial in Nature cited the withdrawal of LYMErix as an instance in which "unfounded public fears place pressures on vaccine developers that go beyond reasonable safety considerations." The original developer of the OspA vaccine at the Max Planck Institute told Nature: "This just shows how irrational the world can be ... There was no scientific justification for the first OspA vaccine LYMErix being pulled." VLA15 The hexavalent (OspA) protein subunit-based vaccine candidate VLA15 was developed by Valneva. It was granted fast track designation by the U.S. Food and Drug Administration in July 2017. In April 2020 Pfizer paid $130 million for the rights to the vaccine, and the companies are developing it together, performing multiple phase 2 trials. A phase 3 trial of VLA15 was scheduled for late 2022, recruiting volunteers at test sites located across the northeastern United States and in Europe. Participants were scheduled to receive an initial three-dose series of vaccines over the course of five to nine months, followed by a booster dose after twelve months, with both the initial series and the booster dose scheduled to be complete before the year's peak Lyme disease season. Other research An mRNA vaccine designed to cause a strong fast immune response to tick saliva allowed the immune system to detect and remove the ticks from test animals before they were able to transmit the infectious bacteria. The vaccine contains mRNAs for the body to build 19 proteins in tick saliva which, by enabling quick development of erythema (itchy redness) at the bite site, protects guinea pigs against Lyme disease. It also protected the test animals if the tick is not removed if only one tick, but not three, remain attached. Canine vaccines Canine vaccines have been formulated and approved for the prevention of Lyme disease in dogs. Currently, three Lyme disease vaccines are available. LymeVax, formulated by Fort Dodge Laboratories, contains intact dead spirochetes which expose the host to the organism. Galaxy Lyme, Intervet-Schering-Plough's vaccine, targets proteins OspC and OspA. The OspC antibodies kill any of the bacteria that have not been killed by the OspA antibodies. Canine Recombinant Lyme, formulated by Merial, generates antibodies against the OspA protein so a tick feeding on a vaccinated dog draws in blood full of anti-OspA antibodies, which kill the spirochetes in the tick's gut before they are transmitted to the dog. Treatment Antibiotics are the primary treatment. The specific approach to their use is dependent on the individual affected and the stage of the disease. For most people with early localized infection, oral administration of doxycycline is widely recommended as the first choice, as it is effective against not only Borrelia bacteria but also a variety of other illnesses carried by ticks. People taking doxycycline should avoid sun exposure because of higher risk of sunburns. Doxycycline is contraindicated in children younger than eight years of age and women who are pregnant or breastfeeding; alternatives to doxycycline are amoxicillin, cefuroxime axetil, and azithromycin. Azithromycin is recommended only in case of intolerance to the other antibiotics. The standard treatment for cellulitis, cephalexin, is not useful for Lyme disease. When it is unclear if a rash is caused by Lyme or cellulitis, the IDSA recommends treatment with cefuroxime or amoxicillin/clavulanic acid, as these are effective against both infections. Individuals with early disseminated or late Lyme infection may have symptomatic cardiac disease, Lyme arthritis, or neurologic symptoms like facial palsy, radiculopathy, meningitis, or peripheral neuropathy. Intravenous administration of ceftriaxone is recommended as the first choice in these cases; cefotaxime and doxycycline are available as alternatives. Treatment regimens for Lyme disease range from 7–14 days in early localized disease, to 14–21 days in early disseminated disease to 14–28 days in late disseminated disease. Neurologic complications of Lyme disease may be treated with doxycycline as it can be taken by mouth and has a lower cost, although in North America evidence of efficacy is only indirect. In case of failure, guidelines recommend retreatment with injectable ceftriaxone. Several months after treatment for Lyme arthritis, if joint swelling persists or returns, a second round of antibiotics may be considered; intravenous antibiotics are preferred for retreatment in case of poor response to oral antibiotics. Outside of that, a prolonged antibiotic regimen lasting more than 28 days is not recommended as no evidence shows it to be effective. IgM and IgG antibody levels may be elevated for years even after successful treatment with antibiotics. As antibody levels are not indicative of treatment success, testing for them is not recommended. Facial palsy may resolve without treatment: however, antibiotic treatment is recommended to stop other Lyme complications. Corticosteroids are not recommended when facial palsy is caused by Lyme disease. In those with facial palsy, frequent use of artificial tears while awake is recommended, along with ointment and a patch or taping the eye closed when sleeping. About a third of people with Lyme carditis need a temporary pacemaker until their heart conduction abnormality resolves, and 21% need to be hospitalized. Lyme carditis should not be treated with corticosteroids. People with Lyme arthritis should limit their level of physical activity to avoid damaging affected joints, and in case of limping should use crutches. Pain associated with Lyme disease may be treated with nonsteroidal anti-inflammatory drugs (NSAIDs). Corticosteroid joint injections are not recommended for Lyme arthritis that is being treated with antibiotics. People with Lyme arthritis treated with intravenous antibiotics or two months of oral antibiotics who continue to have joint swelling two months after treatment and have negative PCR test for Borrelia DNA in the synovial fluid are said to have post-antibiotic Lyme arthritis; this is more common after infection by certain Borrelia strains in people with certain genetic and immunologic characteristics. Post-antibiotic Lyme arthritis may be symptomatically treated with NSAIDs, disease-modifying antirheumatic drugs (DMARDs), arthroscopic synovectomy, or physical therapy. People receiving treatment should be advised that reinfection is possible and how to prevent it. Prognosis Lyme disease's typical first sign, the erythema migrans (EM) rash, resolves within several weeks even without treatment. However, in untreated people, the infection often disseminates to the nervous system, heart or joints, possibly causing permanent damage to body tissues. People who receive recommended antibiotic treatment within several days of appearance of an initial EM rash have the best prospects. Recovery may not be total or immediate. The percentage of people achieving full recovery in the United States increases from about 64–71% at end of treatment for EM rash to about 84–90% after 30 months; higher percentages are reported in Europe. Treatment failure, i.e. persistence of original or appearance of new signs of the disease, occurs only in a few people. Remaining people are considered cured but continue to experience subjective symptoms, e.g. joint or muscle pains or fatigue. These symptoms are usually mild and nondisabling. People treated only after nervous system manifestations of the disease may end up with objective neurological deficits, in addition to subjective symptoms. In Europe, an average of 32–33 months after initial Lyme symptoms in people treated mostly with doxycycline 200 mg for 14–21 days, the percentage of people with lingering symptoms was much higher among those diagnosed with neuroborreliosis (50%) than among those with only an EM rash (16%). In another European study, 5 years after treatment for neuroborreliosis lingering symptoms were less common among children (15%) than adults (30%), and in the latter were less common among those treated within 30 days of the first symptom (16%) than among those treated later (39%); among those with lingering symptoms, 54% had daily activities restricted and 19% were on sick leave or incapacitated. Some data suggest that about 90% of Lyme facial palsies treated with antibiotics recover fully a median of 24 days after appearing and most of the rest recover with only mild abnormality. However, in Europe 41% of people treated for facial palsy had other lingering symptoms at followup up to 6 months later, including 28% with numbness or altered sensation and 14% with fatigue or concentration problems. Palsies in both sides of the face are associated with worse and longer time to recovery. Historical data suggests that untreated people with facial palsies recover at nearly the same rate, but 88% subsequently have Lyme arthritis. Other research shows that synkinesis (involuntary movement of a facial muscle when another one is voluntarily moved) can become evident only 6–12 months after facial palsy appears to be resolved, as damaged nerves regrow and sometimes connect to incorrect muscles. Synkinesis is associated with corticosteroid use. In longer-term follow-up, 16–23% of Lyme facial palsies do not fully recover. In Europe, about a quarter of people with Bannwarth syndrome (Lyme radiculopathy and lymphocytic meningitis) treated with intravenous ceftriaxone for 14 days an average of 30 days after first symptoms had to be retreated 3–6 months later because of unsatisfactory clinical response or continued objective markers of infection in cerebrospinal fluid; after 12 months, 64% recovered fully, 31% had nondisabling mild or infrequent symptoms that did not require regular use of analgesics, and 5% had symptoms that were disabling or required substantial use of analgesics. The most common lingering nondisabling symptoms were headache, fatigue, altered sensation, joint pains, memory disturbances, malaise, radicular pain, sleep disturbances, muscle pains, and concentration disturbances. Lingering disabling symptoms included facial palsy and other impaired movement. Recovery from late neuroborreliosis tends to take longer and be less complete than from early neuroborreliosis, probably because of irreversible neurologic damage. About half the people with Lyme carditis progress to complete heart block, but it usually resolves in a week. Other Lyme heart conduction abnormalities resolve typically within 6 weeks. About 94% of people have full recovery, but 5% need a permanent pacemaker and 1% end up with persistent heart block (the actual percentage may be higher because of unrecognized cases). Lyme myocardial complications usually are mild and self-limiting. However, in some cases Lyme carditis can be fatal. Recommended antibiotic treatments are effective in about 90% of Lyme arthritis cases, although it can take several months for inflammation to resolve and a second round of antibiotics is often necessary. Antibiotic-refractory Lyme arthritis also eventually resolves, typically within 9–14 months (range 4 months – 4 years); DMARDs or synovectomy can accelerate recovery. Reinfection is not uncommon. In a U.S. study, 6–11% of people treated for an EM rash had another EM rash within 30 months. The second rash typically is due to infection by a different Borrelia strain. Post-treatment Lyme disease syndrome Chronic symptoms like pain, fatigue, or cognitive impairment are experienced by 5–20% of people who contract Lyme disease, even after completing treatment. This is called Post-treatment Lyme disease syndrome, or PTLDS. The cause is unknown. One hypothesis is that a persistent, difficult-to-detect infection remains. However, human and animal trials have not provided compelling evidence to support this hypothesis. Another hypothesis is that autoimmunity has been triggered by the infection. Auto–immune responses are known to occur following other infections, including Campylobacter (Guillain-Barré syndrome), Chlamydia (reactive arthritis), and strep throat (rheumatic heart disease). A third hypothesis is that debris from a previous infection could remain. Another hypothesis is that symptoms are simply unrelated to a Lyme infection. Among 13 studies analyzing people who worried about Lyme disease, 47 to 80% had no evidence of Lyme infection while 15 to 55% (median 34%) were able to obtain other diagnoses. There is no proven treatment for Post-treatment Lyme disease syndrome. While short-term antibiotics are effective in early Lyme disease, prolonged antibiotics are not. They have been shown ineffective in placebo-controlled trials and carry the risk of serious, sometimes deadly complications. Generally, treatment is symptomatic and is similar to the management of fibromyalgia or ME/CFS. PTLDS usually gets better over time, but recovery may take many months. Epidemiology Lyme disease occurs regularly in Northern Hemisphere temperate regions. An estimated 476,000 people a year are diagnosed and treated for the disease in the United States. This number is probably an overestimate due to overdiagnosis and overtreatment. Over 200,000 people a year are diagnosed and treated in Europe. There is a suggestion that tick populations and Lyme disease occurrence are increasing and spreading into new areas, owing in part to the warming temperatures of climate change. However, tick-borne disease systems are complex, and determining whether changes are due to climate change or other drivers can be difficult. Lyme disease effects are comparable among males and females. A wide range of age groups is affected, though the number of cases is highest among 10- to 19-year-olds. Africa In northern Africa, B. burgdorferi sensu lato has been identified in Morocco, Algeria, Egypt and Tunisia. Lyme disease in sub-Saharan Africa is presently unknown, but evidence indicates it may occur in humans in this region. The abundance of hosts and tick vectors would favor the establishment of Lyme infection in Africa. In East Africa, two cases of Lyme disease have been reported in Kenya. According The Federation of Infectious Diseases Societies of Southern Africa, Lyme disease is not known to be endemic in either South Africa or Mozambique. Asia B. burgdorferi sensu lato-infested ticks are being found more frequently in Japan, as well as in northwest China, Nepal, Thailand and far eastern Russia. Borrelia has also been isolated in Mongolia. Australia Lyme disease is not considered endemic to Australia. While there have been reports of people acquiring Lyme disease in Australia, and even evidence of closely related Borrelia species in ticks, the evidence linking these cases to local transmission is limited. Ongoing research on resolving potential Borrelia species to Debilitating Symptom Complexes Attributed to Ticks (DSCATT) in Australia are ongoing. Europe In Europe, Lyme disease is caused by infection with one or more pathogenic European genospecies of the spirochaete B. burgdorferi sensu lato, mainly transmitted by the tick Ixodes ricinus. Cases of B. burgdorferi sensu lato-infected ticks are found predominantly in central Europe, particularly in Slovenia and Austria, but have been isolated in almost every country on the continent. The number of cases in southern Europe, such as Italy and Portugal, is much lower. Diagnosed cases in some Western countries, such as Iceland, are rising. Lyme disease is rare in Iceland. On average around 6 to 7 cases are diagnosed every year, primarily localised infections presenting as erythema migrans. None of the cases had a definitive Icelandic origin and the yearly number of cases has not been increasing. United Kingdom In the United Kingdom the number of laboratory-confirmed cases of Lyme disease has been rising steadily since voluntary reporting was introduced in 1986 when 68 cases were recorded in the UK and Ireland combined. In the UK there were 23 confirmed cases in 1988 and 19 in 1990, but 973 in 2009 and 953 in 2010. Provisional figures for the first 3 quarters of 2011 show a 26% increase on the same period in 2010. It is thought, however, that the actual number of cases is significantly higher than suggested by the above figures, with England's Health Protection Agency estimating that there are between 2,000 and 3,000 cases in England and Wales per year (with an average of around 15% of the infections acquired overseas), while Dr Darrel Ho-Yen, Director of the Scottish Toxoplasma Reference Laboratory and National Lyme Disease Testing Service, believes that the number of confirmed cases should be multiplied by 10 "to take account of wrongly diagnosed cases, tests giving false results, sufferers who weren't tested, people who are infected but not showing symptoms, failures to notify and infected individuals who don't consult a doctor." Despite Lyme disease (Borrelia burgdorferi infection) being a notifiable disease in Scotland since January 1990 which should therefore be reported on the basis of clinical suspicion, it is believed that many GPs are unaware of the requirement. Mandatory reporting, limited to laboratory test results only, was introduced throughout the UK in October 2010, under the Health Protection (Notification) Regulations 2010. Although there is a greater number of cases of Lyme disease in the New Forest, Salisbury Plain, Exmoor, the South Downs, parts of Wiltshire and Berkshire, Thetford Forest and the West coast and islands of Scotland, infected ticks are widespread and can even be found in the parks of London. A 1989 report found that 25% of forestry workers in the New Forest were seropositive, as were between 2% and 4–5% of the general local population of the area. Tests on pet dogs carried out throughout the country in 2009 indicated that around 2.5% of ticks in the UK may be infected, considerably higher than previously thought. It is speculated that global warming may lead to an increase in tick activity in the future, as well as an increase in the amount of time that people spend in public parks, thus increasing the risk of infection. However no published research has proven this to be so. North America Many studies in North America have examined ecological and environmental correlates of the number of people affected by Lyme disease. A 2005 study using climate suitability modelling of I. scapularis projected that climate change would cause an overall 213% increase in suitable vector habitat by 2080, with northward expansions in Canada, increased suitability in the central U.S., and decreased suitable habitat and vector retraction in the southern U.S. A 2008 review of published studies concluded that the presence of forests or forested areas was the only variable that consistently elevated the risk of Lyme disease whereas other environmental variables showed little or no concordance between studies. The authors argued that the factors influencing tick density and human risk between sites are still poorly understood, and that future studies should be conducted over longer time periods, become more standardized across regions, and incorporate existing knowledge of regional Lyme disease ecology. Canada The range of ticks able to carry Lyme disease has expanded from a limited area of Ontario to include areas of southern Quebec, Manitoba, northern Ontario, southern New Brunswick, southwest Nova Scotia and limited parts of Saskatchewan and Alberta, as well as British Columbia. Cases have been reported as far east as the island of Newfoundland. A model-based prediction by Leighton et al. (2012) suggests that the range of the I. scapularis tick will expand into Canada by 46 km/year over the next decade, with warming climatic temperatures as the main driver of increased speed of spread. Mexico A 2007 study suggests Borrelia burgdorferi infections are endemic to Mexico, from four cases reported between 1999 and 2000. United States Lyme disease is the most common tick-borne disease in North America and Europe, and one of the fastest-growing infectious diseases in the United States. Of cases reported to the United States CDC, the ratio of Lyme disease infection is 7.9 cases for every 100,000 persons. In the ten states where Lyme disease is most common, the average was 31.6 cases for every 100,000 persons for the year 2005. Although Lyme disease has been reported in all states due to travel-associated infections, about 99% of all reported cases are confined to just five geographic areas (New England, Mid-Atlantic, East-North Central, South Atlantic, and West North-Central). CDC implemented national surveillance of Lyme disease cases in 1991. Since then, reporting criteria has been modified multiple times. The 2022 surveillance case definition classifies cases as confirmed, probable, and suspect. The number of reported cases of the disease has been increasing, as are endemic regions in North America. The CDC emphasizes that, while surveillance data has limitations, it is useful due to "uniformity, simplicity, and timeliness." While cases are under-reported in high-incidence areas, over-reporting is likely in low-incidence areas. Additionally, surveillance cases are reported by county of residence and not where an infection was necessarily contracted. Several similar but apparently distinct conditions may exist, caused by various species or subspecies of Borrelia in North America. A regionally restricted condition that may be related to Borrelia infection is southern tick-associated rash illness (STARI), also known as Masters disease. Amblyomma americanum, known commonly as the lone-star tick, is recognized as the primary vector for STARI. In some parts of the geographical distribution of STARI, Lyme disease is quite rare (e.g., Arkansas), so people in these regions experiencing Lyme-like symptoms—especially if they follow a bite from a lone-star tick—should consider STARI as a possibility. It is generally a milder condition than Lyme and typically responds well to antibiotic treatment. In recent years there have been 5 to 10 cases a year of a disease similar to Lyme occurring in Montana. It occurs primarily in pockets along the Yellowstone River in central Montana. People have developed a red bull's-eye rash around a tick bite followed by weeks of fatigue and a fever. South America In Brazil, Lyme disease is not considered endemic. A Lyme-like disease known as Baggio–Yoshinari syndrome has been described, attributed to microorganisms that do not belong to the B. burgdorferi sensu lato complex and transmitted by ticks of the Amblyomma and Rhipicephalus genera. The first reported case of BYS in Brazil was made in 1992 in Cotia, São Paulo. A 2024 analysis concluded that evidence to connect BYS to Borrelia bacteria was lacking. Etymology Lyme disease was diagnosed as a separate condition for the first time in 1975 in Lyme, Connecticut. History The earliest known evidence of Lyme disease was found in Oetzi, a 5300 year old mummy in the Eastern Alps near the Italian border. The evolutionary history of Borrelia burgdorferi genetics has been examined by scientists. One study has found that prior to the reforestation that accompanied post-colonial farm abandonment in New England and the wholesale migration into the mid-west that occurred during the early 19th century, Lyme disease had been present for thousands of years in America and had spread along with its tick hosts from the Northeast to the Midwest. John Josselyn, who visited New England in 1638 and again from 1663 to 1670, wrote "there be infinite numbers of ticks hanging upon the bushes in summertime that will cleave to man's garments and creep into his breeches, eating themselves in a short time into the very flesh of a man. I have seen the stockings of those that have gone through the woods covered with them." This is also confirmed by the writings of Peter Kalm, a Swedish botanist who was sent to America by Linnaeus, and who found the forests of New York "abound" with ticks when he visited in 1749. When Kalm's journey was retraced 100 years later, the forests were gone and the Lyme bacterium had probably become isolated to a few pockets along the northeast coast, Wisconsin, and Minnesota. Perhaps the first detailed description of what is now known as Lyme disease appeared in the writings of John Walker after a visit to the island of Jura (Deer Island) off the west coast of Scotland in 1764. He gives a good description both of the symptoms of Lyme disease (with "exquisite pain [in] the interior parts of the limbs") and of the tick vector itself, which he describes as a "worm" with a body which is "of a reddish color and of a compressed shape with a row of feet on each side" that "penetrates the skin". Many people from this area of Great Britain emigrated to North America between 1717 and the end of the 18th century. The examination of preserved museum specimens has found Borrelia DNA in an infected Ixodes ricinus tick from Germany that dates back to 1884, and from an infected mouse from Cape Cod that died in 1894. The 2010 autopsy of Ötzi the Iceman, a 5,300-year-old mummy, revealed the presence of the DNA sequence of Borrelia burgdorferi making him the earliest known human with Lyme disease. The early European studies of what is now known as Lyme disease described its skin manifestations. The first study dates to 1883 in Breslau, Germany (now Wrocław, Poland), where physician Alfred Buchwald described a man who for 16 years had had a degenerative skin disorder now known as acrodermatitis chronica atrophicans. At a 1909 research conference, Swedish dermatologist Arvid Afzelius presented a study about an expanding, ring-like lesion he had observed in an older woman following the bite of a sheep tick. He named the lesion erythema migrans. The skin condition now known as borrelial lymphocytoma was first described in 1911. The modern history of medical understanding of the disease, including its cause, diagnosis, and treatment, has been difficult. Neurological problems following tick bites were recognized starting in the 1920s. French physicians Garin and Bujadoux described a farmer with a painful sensory radiculitis accompanied by mild meningitis following a tick bite. A large, ring-shaped rash was also noted, although the doctors did not relate it to the meningoradiculitis. In 1930, the Swedish dermatologist Sven Hellerström was the first to propose EM and neurological symptoms following a tick bite were related. In the 1940s, German neurologist Alfred Bannwarth described several cases of chronic lymphocytic meningitis and polyradiculoneuritis, some of which were accompanied by erythematous skin lesions. Carl Lennhoff, who worked at the Karolinska Institute in Sweden, believed many skin conditions were caused by spirochetes. In 1948, he published on his use of a special stain to microscopically observe what he believed were spirochetes in various types of skin lesions, including EM. Although his conclusions were later shown to be erroneous, interest in the study of spirochetes was sparked. Starting in 1946, facilities in Sweden experimented with treating EM rashes with substances known to kill spirochetes. reported that "penicillin was found to be the most effective." In 1949, Nils Thyresson, who also worked at the Karolinska Institute, was the first to treat ACA with penicillin. In the 1950s, the relationship among tick bite, lymphocytoma, EM and Bannwarth's syndrome was recognized throughout Europe leading to the widespread use of penicillin for treatment in Europe. In 1970 a dermatologist in Wisconsin named Rudolph Scrimenti recognized an EM lesion in a person after recalling a paper by Hellerström that had been reprinted in an American science journal in 1950. This was the first documented case of EM in the United States. Based on the European literature, he treated the person with penicillin. The full syndrome now known as Lyme disease was not recognized until a cluster of cases originally thought to be juvenile rheumatoid arthritis was identified in three towns in southeastern Connecticut in 1975, including the towns Lyme and Old Lyme, which gave the disease its popular name. This was investigated by physicians David Snydman and Allen Steere of the Epidemic Intelligence Service, and by others from Yale University, including Stephen Malawista, who is credited as a co-discover of the disease. The recognition that the people in the United States had EM led to the recognition that "Lyme arthritis" was one manifestation of the same tick-borne condition known in Europe. Before 1976, the elements of B. burgdorferi sensu lato infection were called or known as tick-borne meningopolyneuritis, Garin-Bujadoux syndrome, Bannwarth syndrome, Afzelius's disease, Montauk Knee or sheep tick fever. Since 1976 the disease is most often referred to as Lyme disease, Lyme borreliosis or simply borreliosis. In 1980, Steere, et al., began to test antibiotic regimens in adults with Lyme disease. In the same year, New York State Health Dept. epidemiologist Jorge Benach provided Willy Burgdorfer, a researcher at the Rocky Mountain Biological Laboratory, with collections of I. dammini [scapularis] from Shelter Island, New York, a known Lyme-endemic area as part of an ongoing investigation of Rocky Mountain spotted fever. In examining the ticks for rickettsiae, Burgdorfer noticed "poorly stained, rather long, irregularly coiled spirochetes." Further examination revealed spirochetes in 60% of the ticks. Burgdorfer credited his familiarity with the European literature for his realization that the spirochetes might be the "long-sought cause of ECM and Lyme disease." Benach supplied him with more ticks from Shelter Island and sera from people diagnosed with Lyme disease. University of Texas Health Science Center researcher Alan Barbour "offered his expertise to culture and immunochemically characterize the organism." Burgdorfer subsequently confirmed his discovery by isolating, from people with Lyme disease, spirochetes identical to those found in ticks. In June 1982, he published his findings in Science, and the spirochete was named Borrelia burgdorferi in his honor. After the identification of B. burgdorferi as the causative agent of Lyme disease, antibiotics were selected for testing, guided by in vitro antibiotic sensitivities, including tetracycline antibiotics, amoxicillin, cefuroxime axetil, intravenous and intramuscular penicillin and intravenous ceftriaxone. The mechanism of tick transmission was also the subject of much discussion. B. burgdorferi spirochetes were identified in tick saliva in 1987, confirming the hypothesis that transmission occurred via tick salivary glands. Society, culture, and controversy Landscape changes and urbanization Urbanization and other anthropogenic factors can be implicated in the spread of Lyme disease to humans. In many areas, expansion of suburban neighborhoods has led to gradual deforestation of surrounding wooded areas and increased border contact between humans and tick-dense areas. Human expansion has also resulted in a reduction of predators that hunt deer as well as mice, chipmunks and other small rodents—the primary reservoirs for Lyme disease. As a consequence of increased human contact with host and vector, the likelihood of transmission of the disease has greatly increased. Researchers are investigating possible links between global warming and the spread of vector-borne diseases, including Lyme disease. The dilution effect Given these habitat-host dynamics, in 2003 some researchers began to postulate whether the so called dilution effect could mitigate the spread of Lyme disease. The dilution effect is a hypothesis that predicts that an increase in host biodiversity will result in a decrease in the number of vectors infected with B. burgdorferi. Scientific research has shown that nymphal infection prevalence (NIP) decreases as the number of host species increases, supporting the dilution effect. That said, there is no direct relationship between decreased NIP and decreased epidemiological risk, as this has yet to be proven. Additionally, as of 2018, the dilution effect is only supported in the Northeastern United States, and has been disproved in other parts of the world that also experience high Lyme disease incidence rates Chronic Lyme disease The term "chronic Lyme disease" is controversial and not recognized in the medical literature, and most medical authorities advise against long-term antibiotic treatment for Lyme disease. Studies have shown that most people diagnosed with "chronic Lyme disease" either have no objective evidence of previous or current infection with B. burgdorferi or are people who should be classified as having post-treatment Lyme disease syndrome (PTLDS), which is defined as continuing or relapsing non-specific symptoms (such as fatigue, musculoskeletal pain, and cognitive complaints) in a person previously treated for Lyme disease. The 2008 documentary Under Our Skin is known for promoting controversial and unrecognized theories about "chronic Lyme disease". Other animals Dogs Prevention of Lyme disease is an important step in keeping dogs safe in endemic areas. Prevention education and a number of preventive measures are available. First, for dog owners who live near or who often frequent tick-infested areas, routine vaccinations of their dogs is an important step. Another crucial preventive measure is the use of persistent acaricides, such as topical repellents or pesticides that contain triazapentadienes (Amitraz), phenylpyrazoles (Fipronil), or permethrin (pyrethroids). These acaricides target primarily the adult stages of Lyme-carrying ticks and reduce the number of reproductively active ticks in the environment. Formulations of these ingredients are available in a variety of topical forms, including spot-ons, sprays, powders, impregnated collars, solutions, and shampoos. Examination of a dog for ticks after being in a tick-infested area is an important precautionary measure to take in the prevention of Lyme disease. Key spots to examine include the head, neck, and ears. In dogs, a serious long-term prognosis may result in glomerular disease, which is a category of kidney damage that may cause chronic kidney disease. Dogs may also experience chronic joint disease if the disease is left untreated. However, the majority of cases of Lyme disease in dogs result in complete recovery with, and sometimes without, treatment with antibiotics. In rare cases, Lyme disease can be fatal to both humans and dogs. Cats Unlike dogs, it is very rare for a cat to be infected with Lyme disease. However, cats are nevertheless capable of being infected with B. burgdorferi , following a bite from an infected tick. Cats who are infected with Lyme Disease may show symptoms including but not limited to lameness, fatigue, or loss of appetite. In two cases, the infected cats experienced cardiac irregularities similar to symptoms of Lyme in both dogs and humans. However, cats who are infected with Lyme disease are likely to be asymptomatic, and show no noticeable signs of the disease. Cats with Lyme are often treated with antibiotics, much like other animals. In some cases, additional treatment or therapy may be required. Horses Lyme disease in horses is often challenging to diagnose because symptoms vary widely. Common acute symptoms include weight loss, fever, lameness, ataxia, and other muscle and joint-related issues. Additional symptoms include muscle tenderness, swollen joints, arthritis, and neck stiffness. Chronic symptoms of the disease typically include neurological manifestations, such as meningitis, cranial neuritis, radiculoneuritis, and encephalitis. Furthermore, it is important to note that some horses do not slow clinical signs of Lyme disease. There are three main testing strategies used to diagnose horses with Lyme disease. They include clinical evaluation, serological testing, and polymerase chain reaction (PCR) testing. Detection of specific antibodies against B. burgdorferi alone is not sufficient for a diagnosis of equine Lyme disease and unspecific testing for antibodies is not recommended. Typical treatment involves antibiotics such as oxytetracycline, doxycycline, ceftriaxone, or minocycline. In some cases, a combination of antibiotics may be administered. Doxycycline and minocycline are taken orally, while oxytetracycline and ceftriaxone are usually administered intravenously. The duration and dosage of treatment vary widely among cases. In most cases, the infected horse is euthanized. Death of horses as a result of Borrelia burgdorferi infection remains unknown. Currently, there is no approved Lyme disease vaccine for horses available. However, a study demonstrated that ponies could be protected using an aluminum adjuvanted recombinant outer-surface protein A (rOspA) vaccine. While horses have been administered a Lyme disease vaccine designed for dogs, it elicits only a short-lasting antibody response. Another study supports the use of commercial Lyme disease vaccines, showing that they do elicit an antibody response, which can be significantly enhanced when horses receive an additional booster vaccine.	Biology and health sciences	Infectious disease	null
244362	https://en.wikipedia.org/wiki/Supercell	Supercell	A supercell is a thunderstorm characterized by the presence of a mesocyclone, a deep, persistently rotating updraft. Due to this, these storms are sometimes referred to as rotating thunderstorms. Of the four classifications of thunderstorms (supercell, squall line, multi-cell, and single-cell), supercells are the overall least common and have the potential to be the most severe. Supercells are often isolated from other thunderstorms, and can dominate the local weather up to away. They tend to last 2–4 hours. Supercells are often put into three classification types: classic (normal precipitation level), low-precipitation (LP), and high-precipitation (HP). LP supercells are usually found in climates that are more arid, such as the high plains of the United States, and HP supercells are most often found in moist climates. Supercells can occur anywhere in the world under the right pre-existing weather conditions, but they are most common in the Great Plains of the United States in an area known as Tornado Alley. A high number of supercells are seen in many parts of Europe as well as in the Tornado Corridor (es) of Argentina, Uruguay and southern Brazil. Characteristics Supercells are usually found isolated from other thunderstorms, although they can sometimes be embedded in a squall line. Typically, supercells are found in the warm sector of a low pressure system propagating generally in a north easterly direction in line with the cold front of the low pressure system. Because they can last for hours, they are known as quasi-steady-state storms. Supercells have the capability to deviate from the mean wind. If they track to the right or left of the mean wind (relative to the vertical wind shear), they are said to be "right-movers" or "left-movers," respectively. Supercells can sometimes develop two separate updrafts with opposing rotations, which splits the storm into two supercells: one left-mover and one right-mover. Supercells can be any size – large or small, low or high topped. They usually produce copious amounts of hail, torrential rainfall, strong winds, and substantial downbursts. Supercells are one of the few types of clouds that typically spawn tornadoes within the mesocyclone, although only 30% or fewer do so. Geography Supercells can occur anywhere in the world under the right weather conditions. The first storm to be identified as the supercell type was the Wokingham storm over England, which was studied by Keith Browning and Frank Ludlam in 1962. Browning did the initial work that was followed up by Lemon and Doswell to develop the modern conceptual model of the supercell. To the extent that records are available, supercells are most frequent in the Great Plains of the central United States and southern Canada extending into the southeastern U.S. and northern Mexico; east-central Argentina and adjacent regions of Uruguay; Bangladesh and parts of eastern India; South Africa; and eastern Australia. Supercells occur occasionally in many other mid-latitude regions, including Eastern China and throughout Europe. The areas with highest frequencies of supercells are similar to those with the most occurrences of tornadoes; see tornado climatology and Tornado Alley. Supercell anatomy The current conceptual model of a supercell was described in Severe Thunderstorm Evolution and Mesocyclone Structure as Related to Tornadogenesis by Leslie R. Lemon and Charles A. Doswell III (see Lemon technique). Moisture streams in from the side of the precipitation-free base and merges into a line of warm uplift region where the tower of the thundercloud is tipped by high-altitude shear winds. The high shear causes horizontal vorticity which is tilted within the updraft to become vertical vorticity, and the mass of clouds spins as it gains altitude up to the cap, which can be up to – above ground for the largest storms, and trailing anvil. Supercells derive their rotation through the tilting of horizontal vorticity, which is caused by wind shear imparting rotation upon a rising air parcel by differential forces. Strong updrafts lift the air turning about a horizontal axis and cause this air to turn about a vertical axis. This forms a deep rotating updraft, the mesocyclone. A cap or capping inversion is usually required to form an updraft of sufficient strength. The moisture-laden air is then cooled enough to precipitate as it is rotated toward the cooler region, represented by the turbulent air of the mammatus clouds where the warm air is spilling over top of the cooler, invading air. The cap is formed where shear winds block further uplift for a time, until a relative weakness allows a breakthrough of the cap (an overshooting top); cooler air to the right in the image may or may not form a shelf cloud, but the precipitation zone will occur where the heat engine of the uplift intermingles with the invading, colder air. The cap puts an inverted (warm-above-cold) layer above a normal (cold-above-warm) boundary layer, and by preventing warm surface air from rising, allows one or both of the following: Air below the cap warms and/or becomes more moist Air above the cap cools As the cooler but drier air circulates into the warm, moisture laden inflow, the cloud base will frequently form a wall, and the cloud base often experiences a lowering, which, in extreme cases, are where tornadoes are formed. This creates a warmer, moister layer below a cooler layer, which is increasingly unstable (because warm air is less dense and tends to rise). When the cap weakens or moves, explosive development follows. In North America, supercells usually show up on Doppler weather radar as starting at a point or hook shape on the southwestern side, fanning out to the northeast. The heaviest precipitation is usually on the southwest side, ending abruptly short of the rain-free updraft base or main updraft (not visible to radar). The rear flank downdraft, or RFD, carries precipitation counterclockwise around the north and northwest side of the updraft base, producing a "hook echo" that indicates the presence of a mesocyclone. Structure Overshooting top This "dome" feature appears above the strongest updraft location on the anvil of the storm. It is a result of an updraft powerful enough to break through the upper levels of the troposphere into the lower stratosphere. An observer at ground level and close to the storm may be unable to see the overshooting top because the anvil blocks the sight of this feature. The overshooting is visible from satellite images as a "bubbling" amidst the otherwise smooth upper surface of the anvil cloud. Anvil An anvil forms when the storm's updraft collides with the upper levels of the lowest layer of the atmosphere, or the tropopause, and has nowhere else to go due to the laws of fluid dynamics- specifically pressure, humidity, and density, in simple terms, the packet of air has lost its buoyancy and cannot rise higher. The anvil is very cold (-30°C) and virtually precipitation-free even though virga can be seen falling from the forward sheared anvil. Since there is so little moisture in the anvil, winds can move freely. The clouds take on their anvil shape when the rising air reaches or more. The anvil's distinguishing feature is that it juts out in front of the storm like a shelf. In some cases, it can even shear backwards, called a backsheared anvil, another sign of a very strong updraft. Precipitation-free base This area, typically on the southern side of the storm in North America, is relatively precipitation-free. This is located beneath the main updraft, and is the main area of inflow. While no precipitation may be visible to an observer, large hail may be falling from this area. A region of this area is called the Vault. It is more accurately called the main updraft area. Wall cloud The wall cloud forms near the downdraft/updraft interface. This "interface" is the area between the precipitation area and the precipitation-free base. Wall clouds form when rain-cooled air from the downdraft is pulled into the updraft. This wet, cold air quickly saturates as it is lifted by the updraft, forming a cloud that seems to "descend" from the precipitation-free base. Wall clouds are common and are not exclusive to supercells; only a small percentage actually produce a tornado, but if a storm does produce a tornado, it usually exhibits wall clouds that persist for more than ten minutes. Wall clouds that seem to move violently up or down, and violent movements of cloud fragments (scud or fractus) near the wall cloud, are indications that a tornado could form. Mammatus clouds Mammatus (Mamma, Mammatocumulus) are bulbous or pillow-like cloud formations extending from beneath the anvil of a thunderstorm. These clouds form as cold air in the anvil region of a storm sinks into warmer air beneath it. Mammatus are most apparent when they are lit from one side or below and are therefore at their most impressive near sunset or shortly after sunrise when the sun is low in the sky. Mammatus are not exclusive to supercells and can be associated with developed thunderstorms and cumulonimbus. Forward flank downdraft (FFD) This is generally the area of heaviest and most widespread precipitation. For most supercells, the precipitation core is bounded on its leading edge by a shelf cloud that results from rain-cooled air within the precipitation core spreading outward and interacting with warmer, moist air from outside of the cell. Between the precipitation-free base and the FFD, a "vaulted" or "cathedral" feature can be observed. In high precipitation supercells an area of heavy precipitation may occur beneath the main updraft area where the vault would alternately be observed with classic supercells. Rear flank downdraft (RFD) The rear flank downdraft of a supercell is a very complex and not yet fully understood feature. RFDs mainly occur within classic and HP supercells although RFDs have been observed within LP supercells. The RFD of a supercell is believed to play a large part in tornadogenesis by tightening existing rotation within the surface mesocyclone. RFDs are caused by mid-level steering winds of a supercell colliding with the updraft tower and moving around it in all directions; specifically, the flow that is redirected downward is referred to as the RFD. This downward surge of relatively cool mid-level air, due to interactions between dew points, humidity, and condensation of the converging of air masses, can reach very high speeds and is known to cause widespread wind damage. The radar signature of an RFD is a hook-like structure where sinking air has brought with it precipitation. Flanking line A flanking line is a line of smaller cumulonimbi or cumulus that form in the warm rising air pulled in by the main updraft. Due to convergence and lifting along this line, landspouts sometimes occur on the outflow boundary of this region. Radar features of a supercell Hook echo (or pendant) The "hook echo" is the area of confluence between the main updraft and the rear flank downdraft (RFD). This indicates the position of the mesocyclone and probably a tornado. Bounded weak echo region (or BWER) This is a region of low radar reflectivity bounded above by an area of higher radar reflectivity with an untilted updraft, also called a vault. It is not observed with all supercells but it is at the edge of a very high precipitation echos with a very sharp gradient perpendicular to the RFD. This is evidence of a strong updraft and often the presence of a tornado. To an observer on the ground, it could be experienced as a zone free of precipitation but usually containing large hail. Inflow notch A "notch" of weak reflectivity on the inflow side of the cell. This is not a V-Notch. V Notch A V-shaped notch on the leading edge of the cell, opening away from the main downdraft. This is an indication of divergent flow around a powerful updraft. Hail spike This three body scatter spike is a region of weak echoes found radially behind the main reflectivity core at higher elevations when large hail is present. Descending reflectivity core Supercell variations Supercell thunderstorms are sometimes classified by meteorologists and storm spotters into three categories; however, not all supercells, being hybrid storms, fit neatly into any one category, and many supercells may fall into different categories during different periods of their lifetimes. The standard definition given above is referred to as the Classic supercell. All types of supercells typically produce severe weather. Low precipitation (LP) LP supercells contain a small and relatively light precipitation (rain/hail) core that is well separated from the updraft. The updraft is intense, and LPs are inflow dominant storms. The updraft tower is typically more strongly tilted and the deviant rightward motion less than for other supercell types. The forward flank downdraft (FFD) is noticeably weaker than for other supercell types, and the rear-flank downdraft (RFD) is much weaker—even visually absent in many cases. Like classic supercells, LP supercells tend to form within stronger mid-to-upper level storm-relative wind shear; however, the atmospheric environment leading to their formation is not well understood. The moisture profile of the atmosphere, particularly the depth of the elevated dry layer, also appears to be important, and the low-to-mid level shear may also be important. This type of supercell may be easily identifiable with "sculpted" cloud striations in the updraft base or even a "corkscrewed" or "barber pole" appearance on the updraft, and sometimes an almost "anorexic" look compared to classic supercells. This is because they often form within drier moisture profiles (often initiated by dry lines) leaving LPs with little available moisture despite high mid-to-upper level environmental winds. They most often dissipate rather than turning into classic or HP supercells, although it is still not unusual for LPs to do the latter, especially when moving into a much moister air mass. LPs were first formally described by Howard Bluestein in the early 1980s although storm-chasing scientists noticed them throughout the 1970s. Classic supercells may wither yet maintain updraft rotation as they decay, becoming more like the LP type in a process known as "downscale transition" that also applies to LP storms, and this process is thought to be how many LPs dissipate. LP supercells rarely spawn tornadoes, and those that form tend to be weak, small, and high-based tornadoes, but strong tornadoes have been observed. These storms, although generating lesser precipitation amounts and producing smaller precipitation cores, can generate huge hail. LPs may produce hail larger than baseballs in clear air where no rainfall is visible. LPs are thus hazardous to people and animals caught outside as well as to storm chasers and spotters. Due to the lack of a heavy precipitation core, LP supercells often exhibit relatively weak radar reflectivity without clear evidence of a hook echo, when in fact they are producing a tornado at the time. LP supercells may not even be recognized as supercells in reflectivity data unless one is trained or experienced on their radar characteristics. This is where observations by storm spotter and storm chasers may be of vital importance in addition to Doppler velocity (and polarimetric) radar data. LP supercells are quite sought after by storm chasers because the limited amount of precipitation makes sighting tornadoes at a safe distance much less difficult than with a classic or HP supercell and more so because of the unobscured storm structure unveiled. During spring and early summer, areas in which LP supercells are readily spotted include southwestern Oklahoma and northwestern Texas, among other parts of the western Great Plains. High precipitation (HP) The HP supercell has a much heavier precipitation core that can wrap all the way around the mesocyclone. These are especially dangerous storms, since the mesocyclone is wrapped with rain and can hide a tornado (if present) from view. These storms also cause flooding due to heavy rain, damaging downbursts, and weak tornadoes, although they are also known to produce strong to violent tornadoes. They have a lower potential for damaging hail than Classic and LP supercells, although damaging hail is possible. It has been observed by some spotters that they tend to produce more cloud-to-ground and intracloud lightning than the other types. Also, unlike the LP and Classic types, severe events usually occur at the front (southeast) of the storm. The HP supercell is the most common type of supercell in the United States east of Interstate 35, in the southern parts of the provinces of Ontario and Quebec in Canada, in France, Germany and the Po Valley in north Italy and in the central portions of Argentina and Uruguay. Mini-supercell or low-topped supercell Whereas classic, HP, and LP refer to different precipitation regimes and mesoscale frontal structures, another variation was identified in the early 1990s by Jon Davies. These smaller storms were initially called mini-supercells but are now commonly referred to as low-topped supercells. These are also subdivided into Classic, HP and LP types. Effects Supercells can produce hailstones averaging as large as in diameter, winds over , tornadoes of as strong as EF3 to EF5 intensity (if wind shear and atmospheric instability are able to support the development of stronger tornadoes), flooding, frequent-to-continuous lightning, and very heavy rain. Many tornado outbreaks come from clusters of supercells. Large supercells may spawn multiple long-tracked and deadly tornadoes, with notable examples in the 2011 Super Outbreak. Severe events associated with a supercell almost always occur in the area of the updraft/downdraft interface. In the Northern Hemisphere, this is most often the rear flank (southwest side) of the precipitation area in LP and classic supercells, but sometimes the leading edge (southeast side) of HP supercells. Examples worldwide Asia Some reports suggest that the deluge on 26 July 2005 in Mumbai, India was caused by a supercell when there was a cloud formation high over the city. On this day of rain fell over the city, of which fell in just four hours. The rainfall coincided with a high tide, which exacerbated conditions. Supercells occur commonly from March to May in Bangladesh, West Bengal, and the bordering northeastern Indian states including Tripura. Supercells that produce very high winds with hail and occasional tornadoes are observed in these regions. They also occur along the Northern Plains of India and Pakistan. On March 23, 2013, a massive tornado ripped through Brahmanbaria district in Bangladesh, killing 20 and injuring 200. Australia On New Year's Day 1947 a supercell hit Sydney. The classic type Supercell formed over the Blue Mountains, mid-morning hitting the lower CBD and eastern suburbs by mid-afternoon with the hail similar in size to a cricket ball. At the time, it was the most severe storm to strike the city since recorded observations began in 1792. On April 14, 1999, a severe storm later classified as a supercell hit the east coast of New South Wales. It is estimated that the storm dropped worth of hailstones during its course. At the time it was the most costly disaster in Australia's insurance history, causing an approximated A$2.3 billion worth of damage, of which A$1.7 billion was covered by insurance. On February 27, 2007, a supercell hit Canberra, dumping nearly of ice in Civic. The ice was so heavy that a newly built shopping center's roof collapsed, birds were killed in the hail produced from the supercell, and people were stranded. The following day many homes in Canberra were subjected to flash flooding, caused either by the city's infrastructure's inability to cope with storm water or through mud slides from cleared land. On 6 March 2010, supercell storms hit Melbourne. The storms caused flash flooding in the center of the city and tennis ball-sized () hailstones hit cars and buildings, causing more than $220 million worth of damage and sparking 40,000-plus insurance claims. In just 18 minutes, of rain fell, causing havoc as streets were flooded and trains, planes, and cars were brought to a standstill. That same month, on March 22, 2010 a supercell hit Perth. This storm was one of the worst in the city's history, causing hail stones of in size and torrential rain. The city had its average March rainfall in just seven minutes during the storm. Hail stones caused severe property damage, from dented cars to smashed windows. The storm itself caused more than 100 million dollars in damage. On November 27, 2014 a supercell hit the inner city suburbs including the CBD of Brisbane. Hailstones up to softball size cut power to 71,000 properties, injuring 39 people, and causing a damage bill of $1 billion AUD. A wind gust of 141 km/h was recorded at Archerfield Airport South America An area in South America known as the Tornado Corridor is considered to be the second most frequent location for severe weather, after Tornado Alley in the United States. The region, which covers portions of Argentina, Uruguay, Paraguay, and Brazil during the spring and summer, often experiences strong thunderstorms which may include tornadoes. One of the first known South American supercell thunderstorms to include tornadoes occurred on September 16, 1816, and destroyed the town of Rojas ( west of the city of Buenos Aires). On September 20, 1926, an F4 tornado struck the city of Encarnación (Paraguay), killing over 300 people and making it the second deadliest tornado in South America. On 21 April 1970, the town of Fray Marcos in the Department of Florida, Uruguay experienced an F4 tornado that killed 11, the strongest in the history of the nation. January 10, 1973, saw the most severe tornado in the history of South America: The San Justo tornado, 105 km north of the city of Santa Fe (Argentina), was rated F5, making it the strongest tornado ever recorded in the southern hemisphere, with winds exceeding 400 km/h. On April 13, 1993, in less than 24 hours in the province of Buenos Aires was given the largest tornado outbreak in the history of South America. There were more than 300 tornadoes recorded, with intensities between F1 and F3. The most affected towns were Henderson (EF3), Urdampilleta (EF3) and Mar del Plata (EF2). In December 2000, a series of twelve tornadoes (only registered) affected the Greater Buenos Aires and the province of Buenos Aires, causing serious damage. One of them struck the town of Guernica, and, just two weeks later, in January 2001, an F3 again devastated Guernica, killing 2 people. The December 26, 2003, Tornado F3 happened in Cordoba, with winds exceeding 300 km/h, which hit Córdoba Capital, just 6 km from the city center, in the area known as CPC Route 20, especially neighborhoods of San Roque and Villa Fabric, killing 5 people and injuring hundreds. The tornado that hit the State of São Paulo in 2004 was one of the most destructive in the state, destroying several industrial buildings, 400 houses, killing one and wounding 11. The tornado was rated EF3, but many claim it was a tornado EF4. In November 2009, four tornadoes, rated F1 and F2 reached the town of Posadas (capital of the province of Misiones, Argentina), generating serious damage in the city. Three of the tornadoes affected the airport area, causing damage in Barrio Belén. On April 4, 2012, the Gran Buenos Aires was hit by the storm Buenos Aires, with intensities F1 and F2, which left nearly 30 dead in various locations. On February 21, 2014, in Berazategui (province of Buenos Aires), a tornado of intensity F1 caused material damage including a car was, with two occupants inside, which was elevated a few feet off the ground and flipped over asphalt, both the driver and his passenger were slightly injured. The tornado caused no fatalities. The severe weather that occurred on Tuesday 8/11 had features rarely seen in such magnitude in Argentina. In many towns of La Pampa, San Luis, Buenos Aires and Cordoba, intense hail stones fell up to 6 cm in diameter. On Sunday December 8, 2013, severe storms took place in the center and the coast. The most affected province was Córdoba, storms and supercells type "bow echos" also developed in Santa Fe and San Luis. Europe During the evening of August 3, 2008, a supercell formed over northern France. It spawned an F4 tornado in the Val de Sambre area, about 90 kilometers east of Lille, which impacted nearby cities such as Maubeuge and Hautmont. This same supercell later went on to generate other tornadoes in the Netherlands and Germany. In 2009, on the night of Monday May 25, a supercell formed over Belgium. It was described by Belgian meteorologist Frank Deboosere as "one of the worst storms in recent years" and caused much damage in Belgium – mainly in the provinces of East Flanders (around Ghent), Flemish Brabant (around Brussels) and Antwerp. The storm occurred between about 1:00 am and 4:00 am local time. An incredible 30,000 lightning flashes were recorded in 2 hours – including 10,000 cloud-to-ground strikes. Hailstones up to across were observed in some places and wind gusts over ; in Melle near Ghent a gust of was reported. Trees were uprooted and blown onto several motorways. In Lillo (east of Antwerp) a loaded goods train was blown from the rail tracks. On May 24, 2010, an intense supercell left behind a trail of destruction spanning across three different states in eastern Germany. It produced multiple strong downbursts, damaging hail and at least four tornadoes, most notably an F3 wedge tornado which struck the town of Großenhain, killing one person. On June 28, 2012, three supercells affected England. Two of them formed over the Midlands, producing hailstones reported to be larger than golf balls, with conglomerate stones up to 10 cm across. Burbage in Leicestershire saw some of the most severe hail. Another supercell produced a tornado near Sleaford, in Lincolnshire. On July 28, 2013, an exceptionally long-lived supercell tracked along an almost 400 km long path across parts of Baden-Württemberg and Bavaria in southern Germany, before falling apart in Czechia. The storm had a lifespan of around 7 hours and produced large hail of up to 8 cm in diameter. The city of Reutlingen was hit the hardest, houses and cars were severely damaged, dozens of people injured. With roughly 3.6 billion euros worth of damage, it was by far the costliest thunderstorm event ever documented in Germany. On 25 July 2019 a supercell thunderstorm affected northern England and parts of Northumberland. Large hail, frequent lightning and rotation were reported by many people. On 24 September 2020 a similar event affected parts of West Yorkshire. Only 5 days after that on June 24, 2021, a supercell produced an F4 tornado in south Moravia, Czech Republic. This tornado caused 6 deaths and left more than 200 people injured. With roughly $700 million of damage it was one of the costliest tornadoes to occur outside of the United States. North America Tornado Alley is a region of the central United States where severe weather is common, particularly tornadoes. Supercell thunderstorms occur more frequently in tornado alley and Dixie Alley than anywhere else in the world. Tornado watches and warnings are frequently necessary in the spring and summer. Most places from the Great Plains to the East Coast of the United States and north as far as the Canadian Prairies, the Great Lakes region, and the St. Lawrence River will experience one or more supercells each year. The 1980 Grand Island tornado outbreak affected the city of Grand Island, Nebraska on June 3, 1980. Seven tornadoes touched down in or near the city that night, killing 5 and injuring 200. The Elie, Manitoba tornado was an F5 that struck the town of Elie, Manitoba on June 22, 2007. While several houses were leveled, no one was injured or killed by the tornado. The most intense tornado outbreaks on record, known as super outbreaks, have all occurred in the United States. The 1974 Super Outbreak and 2011 Super Outbreak each spawned over 10 violent tornadoes, killed over 300, and caused billions in damage, most of which can be attributed to tornado damage. A massive tornado outbreak on May 3, 1999 spawned an F5 tornado in the area of Oklahoma City that had the highest recorded winds on Earth. Another series of tornadoes, which occurred in May 2013, caused severe devastation to Oklahoma City in general. From May 18 to May 21, a series of tornadoes hit, including a tornado which was later rated EF5, which traveled across parts of the Oklahoma City area, causing a severe amount of damage in a heavily populated section of Moore. Twenty-three fatalities and 377 injuries were caused by the tornado. Sixty-one other tornadoes were confirmed during the storm period. Later on in the same month, on the night of May 31, 2013, another eight deaths were confirmed from what became the widest tornado on record which hit El Reno, Oklahoma, one of a series of tornadoes and funnel clouds which hit nearby areas. In Mexico, the tallest non-tropical thunderstorm on record occurred as a high-topped supercell near Nueva Rosita, Coahuila on May 24, 2016. This storm was recorded at a height of and produced lightning as far away as from the center. South Africa South Africa witnesses several supercell thunderstorms each year with the inclusion of isolated tornadoes. On most occasions these tornadoes occur in open farmlands and rarely cause damage to property, as such many of the tornadoes which do occur in South Africa are not reported. The majority of supercells develop in the central, northern, and north eastern parts of the country. The Free State, Gauteng, and Kwazulu Natal are typically the provinces where these storms are most commonly experienced, though supercell activity is not limited to these provinces. On occasion, hail reaches sizes in excess of golf balls, and tornadoes, though rare, also occur. On 6 May 2009, a well-defined hook echo was noticed on local South African radars, along with satellite imagery this supported the presence of a strong supercell storm. Reports from the area indicated heavy rains, winds and large hail. On October 2, 2011, two devastating tornadoes tore through two separate parts of South Africa on the same day, hours apart from each other. The first, classified as an EF2 hit Meqheleng, the informal settlement outside Ficksburg, Free State which devastated shacks and homes, uprooted trees, and killed one small child. The second, which hit the informal settlement of Duduza, Nigel in the Gauteng province, also classified as EF2 hit hours apart from the one that struck Ficksburg. This tornado completely devastated parts of the informal settlement and killed two children, destroying shacks and RDP homes. Gallery	Physical sciences	Storms	Earth science
244373	https://en.wikipedia.org/wiki/Scalar%20multiplication	Scalar multiplication	In mathematics, scalar multiplication is one of the basic operations defining a vector space in linear algebra (or more generally, a module in abstract algebra). In common geometrical contexts, scalar multiplication of a real Euclidean vector by a positive real number multiplies the magnitude of the vector without changing its direction. Scalar multiplication is the multiplication of a vector by a scalar (where the product is a vector), and is to be distinguished from inner product of two vectors (where the product is a scalar). Definition In general, if K is a field and V is a vector space over K, then scalar multiplication is a function from K × V to V. The result of applying this function to k in K and v in V is denoted kv. Properties Scalar multiplication obeys the following rules (vector in boldface): Additivity in the scalar: (c + d)v = cv + dv; Additivity in the vector: c(v + w) = cv + cw; Compatibility of product of scalars with scalar multiplication: (cd)v = c(dv); Multiplying by 1 does not change a vector: 1v = v; Multiplying by 0 gives the zero vector: 0v = 0; Multiplying by −1 gives the additive inverse: (−1)v = −v. Here, + is addition either in the field or in the vector space, as appropriate; and 0 is the additive identity in either. Juxtaposition indicates either scalar multiplication or the multiplication operation in the field. Interpretation The space of vectors may be considered a coordinate space where elements are associated with a list of elements from K. The units of the field form a group K × and the scalar-vector multiplication is a group action on the coordinate space by K ×. The zero of the field acts on the coordinate space to collapse it to the zero vector. When K is the field of real numbers there is a geometric interpretation of scalar multiplication: it stretches or contracts vectors by a constant factor. As a result, it produces a vector in the same or opposite direction of the original vector but of a different length. As a special case, V may be taken to be K itself and scalar multiplication may then be taken to be simply the multiplication in the field. When V is Kn, scalar multiplication is equivalent to multiplication of each component with the scalar, and may be defined as such. The same idea applies if K is a commutative ring and V is a module over K. K can even be a rig, but then there is no additive inverse. If K is not commutative, the distinct operations left scalar multiplication cv and right scalar multiplication vc may be defined. Scalar multiplication of matrices The left scalar multiplication of a matrix with a scalar gives another matrix of the same size as . It is denoted by , whose entries of are defined by explicitly: Similarly, even though there is no widely-accepted definition, the right scalar multiplication of a matrix with a scalar could be defined to be explicitly: When the entries of the matrix and the scalars are from the same commutative field, for example, the real number field or the complex number field, these two multiplications are the same, and can be simply called scalar multiplication. For matrices over a more general field that is not commutative, they may not be equal. For a real scalar and matrix: For quaternion scalars and matrices: where are the quaternion units. The non-commutativity of quaternion multiplication prevents the transition of changing to .	Mathematics	Linear algebra	null
244374	https://en.wikipedia.org/wiki/X86-64	X86-64	x86-64 (also known as x64, x86_64, AMD64, and Intel 64) is a 64-bit extension of the x86 instruction set architecture first announced in 1999. It introduces two new operating modes: 64-bit mode and compatibility mode, along with a new four-level paging mechanism. In 64-bit mode, x86-64 supports significantly large amounts of virtual memory and physical memory compared to its 32-bit predecessors, allowing programs to utilize more memory for data storage. The architecture expands the number of general-purpose registers from 8 to 16, all fully general-purpose, and extends their width to 64 bits. Floating-point arithmetic is supported through mandatory SSE2 instructions in 64-bit mode. While the older x87 FPU and MMX registers are still available, they are generally superceded by a set of sixteen 128-bit vector registers (XMM registers). Each of these vector registers can store one or two double-precision floating-point numbers, up to four single-precision floating-point numbers, or various integer formats. In 64-bit mode, instructions are modified to support 64-bit operands and 64-bit addressing mode. The x86-64 architecture defines a compatibility mode that allows 16-bit and 32-bit user applications to run unmodified alongside 64-bit applications, provided the 64-bit operating system supports them. Since the full x86-32 instruction sets remain implemented in hardware without the need for emulation, these older executables can run with little or no performance penalty, while newer or modified applications can take advantage of new features of the processor design to achieve performance improvements. Also, processors supporting x86-64 still power on in real mode to maintain backward compatibility with the original 8086 processor, as has been the case with x86 processors since the introduction of protected mode with the 80286. The original specification, created by AMD and released in 2000, has been implemented by AMD, Intel, and VIA. The AMD K8 microarchitecture, in the Opteron and Athlon 64 processors, was the first to implement it. This was the first significant addition to the x86 architecture designed by a company other than Intel. Intel was forced to follow suit and introduced a modified NetBurst family which was software-compatible with AMD's specification. VIA Technologies introduced x86-64 in their VIA Isaiah architecture, with the VIA Nano. The x86-64 architecture was quickly adopted for desktop and laptop personal computers and servers which were commonly configured for 16 GiB (gibibytes) of memory or more. It has effectively replaced the discontinued Intel Itanium architecture (formerly IA-64), which was originally intended to replace the x86 architecture. x86-64 and Itanium are not compatible on the native instruction set level, and operating systems and applications compiled for one architecture cannot be run on the other natively. AMD64 History AMD64 (also variously referred to by AMD in their literature and documentation as “AMD 64-bit Technology” and “AMD x86-64 Architecture”) was created as an alternative to the radically different IA-64 architecture designed by Intel and Hewlett-Packard, which was backward-incompatible with IA-32, the 32-bit version of the x86 architecture. AMD originally announced AMD64 in 1999 with a full specification available in August 2000. As AMD was never invited to be a contributing party for the IA-64 architecture and any kind of licensing seemed unlikely, the AMD64 architecture was positioned by AMD from the beginning as an evolutionary way to add 64-bit computing capabilities to the existing x86 architecture while supporting legacy 32-bit x86 code, as opposed to Intel's approach of creating an entirely new, completely x86-incompatible 64-bit architecture with IA-64. The first AMD64-based processor, the Opteron, was released in April 2003. Implementations AMD's processors implementing the AMD64 architecture include Opteron, Athlon 64, Athlon 64 X2, Athlon 64 FX, Athlon II (followed by "X2", "X3", or "X4" to indicate the number of cores, and XLT models), Turion 64, Turion 64 X2, Sempron ("Palermo" E6 stepping and all "Manila" models), Phenom (followed by "X3" or "X4" to indicate the number of cores), Phenom II (followed by "X2", "X3", "X4" or "X6" to indicate the number of cores), FX, Fusion/APU and Ryzen/Epyc. Architectural features The primary defining characteristic of AMD64 is the availability of 64-bit general-purpose processor registers (for example, ), 64-bit integer arithmetic and logical operations, and 64-bit virtual addresses. The designers took the opportunity to make other improvements as well. Notable changes in the 64-bit extensions include: 64-bit integer capability All general-purpose registers (GPRs) are expanded from 32 bits to 64 bits, and all arithmetic and logical operations, memory-to-register and register-to-memory operations, etc., can operate directly on 64-bit integers. Pushes and pops on the stack default to 8-byte strides, and pointers are 8 bytes wide. Additional registers In addition to increasing the size of the general-purpose registers, the number of named general-purpose registers is increased from eight (i.e. , , , , , , , ) in x86 to 16 (i.e. , , , , , , , , , , , , , , , ). It is therefore possible to keep more local variables in registers rather than on the stack, and to let registers hold frequently accessed constants; arguments for small and fast subroutines may also be passed in registers to a greater extent. AMD64 still has fewer registers than many RISC instruction sets (e.g. Power ISA has 32 GPRs; 64-bit ARM, RISC-V I, SPARC, Alpha, MIPS, and PA-RISC have 31) or VLIW-like machines such as the IA-64 (which has 128 registers). However, an AMD64 implementation may have far more internal registers than the number of architectural registers exposed by the instruction set (see register renaming). (For example, AMD Zen cores have 168 64-bit integer and 160 128-bit vector floating-point physical internal registers.) Additional XMM (SSE) registers Similarly, the number of 128-bit XMM registers (used for Streaming SIMD instructions) is also increased from 8 to 16. The traditional x87 FPU register stack is not included in the register file size extension in 64-bit mode, compared with the XMM registers used by SSE2, which did get extended. The x87 register stack is not a simple register file although it does allow direct access to individual registers by low cost exchange operations. Larger virtual address space The AMD64 architecture defines a 64-bit virtual address format, of which the low-order 48 bits are used in current implementations. This allows up to 256 TiB (248 bytes) of virtual address space. The architecture definition allows this limit to be raised in future implementations to the full 64 bits, extending the virtual address space to 16 EiB (264 bytes). This is compared to just 4 GiB (232 bytes) for the x86. This means that very large files can be operated on by mapping the entire file into the process's address space (which is often much faster than working with file read/write calls), rather than having to map regions of the file into and out of the address space. Larger physical address space The original implementation of the AMD64 architecture implemented 40-bit physical addresses and so could address up to 1 TiB (240 bytes) of RAM. Current implementations of the AMD64 architecture (starting from AMD 10h microarchitecture) extend this to 48-bit physical addresses and therefore can address up to 256 TiB (248 bytes) of RAM. The architecture permits extending this to 52 bits in the future (limited by the page table entry format); this would allow addressing of up to 4 PiB of RAM. For comparison, 32-bit x86 processors are limited to 64 GiB of RAM in Physical Address Extension (PAE) mode, or 4 GiB of RAM without PAE mode. Larger physical address space in legacy mode When operating in legacy mode the AMD64 architecture supports Physical Address Extension (PAE) mode, as do most current x86 processors, but AMD64 extends PAE from 36 bits to an architectural limit of 52 bits of physical address. Any implementation, therefore, allows the same physical address limit as under long mode. Instruction pointer relative data access Instructions can now reference data relative to the instruction pointer (RIP register). This makes position-independent code, as is often used in shared libraries and code loaded at run time, more efficient. SSE instructions The original AMD64 architecture adopted Intel's SSE and SSE2 as core instructions. These instruction sets provide a vector supplement to the scalar x87 FPU, for the single-precision and double-precision data types. SSE2 also offers integer vector operations, for data types ranging from 8bit to 64bit precision. This makes the vector capabilities of the architecture on par with those of the most advanced x86 processors of its time. These instructions can also be used in 32-bit mode. The proliferation of 64-bit processors has made these vector capabilities ubiquitous in home computers, allowing the improvement of the standards of 32-bit applications. The 32-bit edition of Windows 8, for example, requires the presence of SSE2 instructions. SSE3 instructions and later Streaming SIMD Extensions instruction sets are not standard features of the architecture. No-Execute bit The No-Execute bit or NX bit (bit 63 of the page table entry) allows the operating system to specify which pages of virtual address space can contain executable code and which cannot. An attempt to execute code from a page tagged "no execute" will result in a memory access violation, similar to an attempt to write to a read-only page. This should make it more difficult for malicious code to take control of the system via "buffer overrun" or "unchecked buffer" attacks. A similar feature has been available on x86 processors since the 80286 as an attribute of segment descriptors; however, this works only on an entire segment at a time. Segmented addressing has long been considered an obsolete mode of operation, and all current PC operating systems in effect bypass it, setting all segments to a base address of zero and (in their 32-bit implementation) a size of 4 GiB. AMD was the first x86-family vendor to implement no-execute in linear addressing mode. The feature is also available in legacy mode on AMD64 processors, and recent Intel x86 processors, when PAE is used. Removal of older features A few "system programming" features of the x86 architecture were either unused or underused in modern operating systems and are either not available on AMD64 in long (64-bit and compatibility) mode, or exist only in limited form. These include segmented addressing (although the FS and GS segments are retained in vestigial form for use as extra-base pointers to operating system structures), the task state switch mechanism, and virtual 8086 mode. These features remain fully implemented in "legacy mode", allowing these processors to run 32-bit and 16-bit operating systems without modifications. Some instructions that proved to be rarely useful are not supported in 64-bit mode, including saving/restoring of segment registers on the stack, saving/restoring of all registers (PUSHA/POPA), decimal arithmetic, BOUND and INTO instructions, and "far" jumps and calls with immediate operands. Virtual address space details Canonical form addresses Although virtual addresses are 64 bits wide in 64-bit mode, current implementations (and all chips that are known to be in the planning stages) do not allow the entire virtual address space of 264 bytes (16 EiB) to be used. This would be approximately four billion times the size of the virtual address space on 32-bit machines. Most operating systems and applications will not need such a large address space for the foreseeable future, so implementing such wide virtual addresses would simply increase the complexity and cost of address translation with no real benefit. AMD, therefore, decided that, in the first implementations of the architecture, only the least significant 48 bits of a virtual address would actually be used in address translation (page table lookup). In addition, the AMD specification requires that the most significant 16 bits of any virtual address, bits 48 through 63, must be copies of bit 47 (in a manner akin to sign extension). If this requirement is not met, the processor will raise an exception. Addresses complying with this rule are referred to as "canonical form." Canonical form addresses run from 0 through 00007FFF'FFFFFFFF, and from FFFF8000'00000000 through FFFFFFFF'FFFFFFFF, for a total of 256 TiB of usable virtual address space. This is still 65,536 times larger than the virtual 4 GiB address space of 32-bit machines. This feature eases later scalability to true 64-bit addressing. Many operating systems (including, but not limited to, the Windows NT family) take the higher-addressed half of the address space (named kernel space) for themselves and leave the lower-addressed half (user space) for application code, user mode stacks, heaps, and other data regions. The "canonical address" design ensures that every AMD64 compliant implementation has, in effect, two memory halves: the lower half starts at 00000000'00000000 and "grows upwards" as more virtual address bits become available, while the higher half is "docked" to the top of the address space and grows downwards. Also, enforcing the "canonical form" of addresses by checking the unused address bits prevents their use by the operating system in tagged pointers as flags, privilege markers, etc., as such use could become problematic when the architecture is extended to implement more virtual address bits. The first versions of Windows for x64 did not even use the full 256 TiB; they were restricted to just 8 TiB of user space and 8 TiB of kernel space. Windows did not support the entire 48-bit address space until Windows 8.1, which was released in October 2013. Page table structure The 64-bit addressing mode ("long mode") is a superset of Physical Address Extensions (PAE); because of this, page sizes may be 4 KiB (212 bytes) or 2 MiB (221 bytes). Long mode also supports page sizes of 1 GiB (230 bytes). Rather than the three-level page table system used by systems in PAE mode, systems running in long mode use four levels of page table: PAE's Page-Directory Pointer Table is extended from four entries to 512, and an additional Page-Map Level 4 (PML4) Table is added, containing 512 entries in 48-bit implementations. A full mapping hierarchy of 4 KiB pages for the whole 48-bit space would take a bit more than 512 GiB of memory (about 0.195% of the 256 TiB virtual space). {\| class="wikitable" style="text-align:center" \|+ style="font-size: 105%; white-space: nowrap;" \| 64 bit page table entry \|- style="border-top: 2px solid #777777;" ! Bits: ! 63 ! colspan="11" \| 62 … 52 ! colspan="20" \| 51 … 32 \|- style="border-bottom: 2px solid #777777;" ! Content: \| NX \| colspan="11" style="background-color: #ccc;" \| reserved \| colspan="20" \| Bit 51…32 of base address \|- !Bits: ! colspan="20" \| 31 … 12 ! colspan="3" \| 11 … 9 ! 8 ! 7 ! 6 ! 5 ! 4 ! 3 ! 2 ! 1 ! 0 \|- style="border-bottom: 2px solid #777777;" ! Content: \| colspan="20" \| Bit 31…12 of base address \| colspan="3" \| ign. \| G \| PAT \| D \| A \| PCD \| PWT \| U/S \| R/W \| P \|} Intel has implemented a scheme with a 5-level page table, which allows Intel 64 processors to support 57-bit addresses, and in turn, a 128 PiB virtual address space. Further extensions may allow full 64-bit virtual address space and physical memory with 12-bit page table descriptors and 16- or 21-bit memory offsets for 64 KiB and 2 MiB page allocation sizes; the page table entry would be expanded to 128 bits to support additional hardware flags for page size and virtual address space size. Operating system limits The operating system can also limit the virtual address space. Details, where applicable, are given in the "Operating system compatibility and characteristics" section. Physical address space details Current AMD64 processors support a physical address space of up to 248 bytes of RAM, or 256 TiB. However, , there were no known x86-64 motherboards that support 256 TiB of RAM. The operating system may place additional limits on the amount of RAM that is usable or supported. Details on this point are given in the "Operating system compatibility and characteristics" section of this article. Operating modes The architecture has two primary modes of operation: long mode and legacy mode. Long mode Long mode is the architecture's intended primary mode of operation; it is a combination of the processor's native 64-bit mode and a combined 32-bit and 16-bit compatibility mode. It is used by 64-bit operating systems. Under a 64-bit operating system, 64-bit programs run under 64-bit mode, and 32-bit and 16-bit protected mode applications (that do not need to use either real mode or virtual 8086 mode in order to execute at any time) run under compatibility mode. Real-mode programs and programs that use virtual 8086 mode at any time cannot be run in long mode unless those modes are emulated in software. However, such programs may be started from an operating system running in long mode on processors supporting VT-x or AMD-V by creating a virtual processor running in the desired mode. Since the basic instruction set is the same, there is almost no performance penalty for executing protected mode x86 code. This is unlike Intel's IA-64, where differences in the underlying instruction set mean that running 32-bit code must be done either in emulation of x86 (making the process slower) or with a dedicated x86 coprocessor. However, on the x86-64 platform, many x86 applications could benefit from a 64-bit recompile, due to the additional registers in 64-bit code and guaranteed SSE2-based FPU support, which a compiler can use for optimization. However, applications that regularly handle integers wider than 32 bits, such as cryptographic algorithms, will need a rewrite of the code handling the huge integers in order to take advantage of the 64-bit registers. Legacy mode Legacy mode is the mode that the processor is in when it is not in long mode. In this mode, the processor acts like an older x86 processor, and only 16-bit and 32-bit code can be executed. Legacy mode allows for a maximum of 32 bit virtual addressing which limits the virtual address space to 4 GiB. 64-bit programs cannot be run from legacy mode. Protected mode Protected mode is made into a submode of legacy mode. It is the submode that 32-bit operating systems and 16-bit protected mode operating systems operate in when running on an x86-64 CPU. Real mode Real mode is the initial mode of operation when the processor is initialized, and is a submode of legacy mode. It is backwards compatible with the original Intel 8086 and Intel 8088 processors. Real mode is primarily used today by operating system bootloaders, which are required by the architecture to configure virtual memory details before transitioning to higher modes. This mode is also used by any operating system that needs to communicate with the system firmware with a traditional BIOS-style interface. Intel 64 Intel 64 is Intel's implementation of x86-64, used and implemented in various processors made by Intel. History Historically, AMD has developed and produced processors with instruction sets patterned after Intel's original designs, but with x86-64, roles were reversed: Intel found itself in the position of adopting the ISA that AMD created as an extension to Intel's own x86 processor line. Intel's project was originally codenamed Yamhill (after the Yamhill River in Oregon's Willamette Valley). After several years of denying its existence, Intel announced at the February 2004 IDF that the project was indeed underway. Intel's chairman at the time, Craig Barrett, admitted that this was one of their worst-kept secrets. Intel's name for this instruction set has changed several times. The name used at the IDF was CT (presumably for Clackamas Technology, another codename from an Oregon river); within weeks they began referring to it as IA-32e (for IA-32 extensions) and in March 2004 unveiled the "official" name EM64T (Extended Memory 64 Technology). In late 2006 Intel began instead using the name Intel 64 for its implementation, paralleling AMD's use of the name AMD64. The first processor to implement Intel 64 was the multi-socket processor Xeon code-named Nocona in June 2004. In contrast, the initial Prescott chips (February 2004) did not enable this feature. Intel subsequently began selling Intel 64-enabled Pentium 4s using the E0 revision of the Prescott core, being sold on the OEM market as the Pentium 4, model F. The E0 revision also adds eXecute Disable (XD) (Intel's name for the NX bit) to Intel 64, and has been included in then current Xeon code-named Irwindale. Intel's official launch of Intel 64 (under the name EM64T at that time) in mainstream desktop processors was the N0 stepping Prescott-2M. The first Intel mobile processor implementing Intel 64 is the Merom version of the Core 2 processor, which was released on July 27, 2006. None of Intel's earlier notebook CPUs (Core Duo, Pentium M, Celeron M, Mobile Pentium 4) implement Intel 64. Implementations Intel's processors implementing the Intel64 architecture include the Pentium 4 F-series/5x1 series, 506, and 516, Celeron D models 3x1, 3x6, 355, 347, 352, 360, and 365 and all later Celerons, all models of Xeon since "Nocona", all models of Pentium Dual-Core processors since "Merom-2M", the Atom 230, 330, D410, D425, D510, D525, N450, N455, N470, N475, N550, N570, N2600 and N2800, all versions of the Pentium D, Pentium Extreme Edition, Core 2, Core i9, Core i7, Core i5, and Core i3 processors, and the Xeon Phi 7200 series processors. X86S X86S was a simplification of x86-64 first proposed by Intel in May 2023. The new architecture would have removed support for 16-bit and 32-bit operating systems, although 32-bit programs would still run under a 64-bit OS. A compliant CPU would have no longer had legacy mode, and started directly in 64-bit long mode. There would have been a way to switch to 5-level paging without going through the unpaged mode. Specific removed features included: Segmentation gates 32-bit ring 0 VT-x will no longer emulate this feature Rings 1 and 2 Ring 3 I/O port (/) access; see port-mapped I/O String port I/O (/) Real mode (including huge real mode), 16-bit protected mode, VM86 16-bit addressing mode VT-x will no longer provide unrestricted mode 8259 support; the only APIC supported would be X2APIC Some unused operating system mode bits 16-bit and 32-bit Startup IPI (SIPI) The draft specification received multiple updates, reaching version 1.2 by June 2024. It was eventually abandoned as of December 2024, following the formation of the x86 Ecosystem Advisory Group by Intel and AMD. Advanced Performance Extensions Advanced Performance Extensions is a 2023 Intel proposal for new instructions and an additional 16 general-purpose registers. VIA's x86-64 implementation VIA Technologies introduced their first implementation of the x86-64 architecture in 2008 after five years of development by its CPU division, Centaur Technology. Codenamed "Isaiah", the 64-bit architecture was unveiled on January 24, 2008, and launched on May 29 under the VIA Nano brand name. The processor supports a number of VIA-specific x86 extensions designed to boost efficiency in low-power appliances. It is expected that the Isaiah architecture will be twice as fast in integer performance and four times as fast in floating-point performance as the previous-generation VIA Esther at an equivalent clock speed. Power consumption is also expected to be on par with the previous-generation VIA CPUs, with thermal design power ranging from 5 W to 25 W. Being a completely new design, the Isaiah architecture was built with support for features like the x86-64 instruction set and x86 virtualization which were unavailable on its predecessors, the VIA C7 line, while retaining their encryption extensions. Microarchitecture levels In 2020, through a collaboration between AMD, Intel, Red Hat, and SUSE, three microarchitecture levels (or feature levels) on top of the x86-64 baseline were defined: x86-64-v2, x86-64-v3, and x86-64-v4. These levels define specific features that can be targeted by programmers to provide compile-time optimizations. The features exposed by each level are as follows: The x86-64 microarchitecture feature levels can also be found as AMD64-v1, AMD64-v2 .. or AMD64_v1 .. in settings where the "AMD64" nomenclature is used. These are used as synonyms with the x86-64-vX nomenclature and are thus functionally identical. E.g. the Go language documentation or the Fedora linux distribution. All levels include features found in the previous levels. Instruction set extensions not concerned with general-purpose computation, including AES-NI and RDRAND, are excluded from the level requirements. Differences between AMD64 and Intel 64 Although nearly identical, there are some differences between the two instruction sets in the semantics of a few seldom used machine instructions (or situations), which are mainly used for system programming. Compilers generally produce executables (i.e. machine code) that avoid any differences, at least for ordinary application programs. This is therefore of interest mainly to developers of compilers, operating systems and similar, which must deal with individual and special system instructions. Recent implementations Intel 64 allows SYSCALL/SYSRET only in 64-bit mode (not in compatibility mode), and allows SYSENTER/SYSEXIT in both modes. AMD64 lacks SYSENTER/SYSEXIT in both sub-modes of long mode. When returning to a non-canonical address using SYSRET, AMD64 processors execute the general protection fault handler in privilege level 3, while on Intel 64 processors it is executed in privilege level 0. AMD64 requires a different microcode update format and control MSRs (model-specific registers), while Intel 64 implements microcode update unchanged from their 32-bit only processors. Intel 64 lacks some MSRs that are considered architectural in AMD64. These include SYSCFG, TOP_MEM, and TOP_MEM2. Intel 64 lacks the ability to save and restore a reduced (and thus faster) version of the floating-point state (involving the FXSAVE and FXRSTOR instructions). In 64-bit mode, near branches with the 66H (operand size override) prefix behave differently. Intel 64 ignores this prefix: the instruction has a 32-bit sign extended offset, and instruction pointer is not truncated. AMD64 uses a 16-bit offset field in the instruction, and clears the top 48 bits of instruction pointer. On Intel 64 but not AMD64, the REX.W prefix can be used with the far-pointer instructions (LFS, LGS, LSS, , ) to increase the size of their far pointer argument to 80 bits (64-bit offset + 16-bit segment). When the MOVSXD instruction is executed with a memory source operand and an operand-size of 16 bits, the memory operand will be accessed with a 16-bit read on Intel 64, but a 32-bit read on AMD64. The FCOMI/FCOMIP/FUCOMI/FUCOMIP (x87 floating-point compare) instructions will clear the OF, SF and AF bits of EFLAGS on Intel 64, but leave these flag bits unmodified on AMD64. For the VMASKMOVPS/VMASKMOVPD/VPMASKMOVD/VPMASKMOVQ (AVX/AVX2 masked move to/from memory) instructions, Intel 64 architecturally guarantees that the instructions will not cause memory faults (e.g. page-faults and segmentation-faults) for any zero-masked lanes, while AMD64 does not provide such a guarantee. If the RDRAND instruction fails to obtain a random number (as indicated by EFLAGS.CF=0), the destination register is architecturally guaranteed to be set to 0 on Intel 64 but not AMD64. For the VPINSRD and VPEXTRD (AVX vector lane insert/extract) instructions outside 64-bit mode, AMD64 requires the instructions to be encoded with VEX.W=0, while Intel 64 also accepts encodings with VEX.W=1. (In 64-bit mode, both AMD64 and Intel 64 require VEX.W=0.) The 0F 0D /r opcode with the ModR/M byte's Mod field set to 11b is a Reserved-NOP on Intel 64 but will cause #UD (invalid-opcode exception) on AMD64. The ordering guarantees provided by some memory ordering instructions such as LFENCE and MFENCE differ between Intel 64 and AMD64: LFENCE is dispatch-serializing (enabling it to be used as a speculation fence) on Intel 64 but is not architecturally guaranteed to be dispatch-serializing on AMD64. MFENCE is a fully serializing instruction (including instruction fetch serialization) on AMD64 but not Intel 64. The MOV to CR8 and INVPCID instructions are serializing on AMD64 but not Intel 64. The LMSW instruction is serializing on Intel 64 but not AMD64. Older implementations The AMD64 processors prior to Revision F (distinguished by the switch from DDR to DDR2 memory and new sockets AM2, F and S1) of 2006 lacked the CMPXCHG16B instruction, which is an extension of the CMPXCHG8B instruction present on most post-80486 processors. Similar to CMPXCHG8B, CMPXCHG16B allows for atomic operations on octa-words (128-bit values). This is useful for parallel algorithms that use compare and swap on data larger than the size of a pointer, common in lock-free and wait-free algorithms. Without CMPXCHG16B one must use workarounds, such as a critical section or alternative lock-free approaches. Its absence also prevents 64-bit Windows prior to Windows 8.1 from having a user-mode address space larger than 8 TiB. The 64-bit version of Windows 8.1 requires the instruction. Early AMD64 and Intel 64 CPUs lacked LAHF and SAHF instructions in 64-bit mode. AMD introduced these instructions (also in 64-bit mode) with their 90 nm (revision D) processors, starting with Athlon 64 in October 2004. Intel introduced the instructions in October 2005 with the 0F47h and later revisions of NetBurst. The 64-bit version of Windows 8.1 requires this feature. Early Intel CPUs with Intel 64 also lack the NX bit of the AMD64 architecture. It was added in the stepping E0 (0F41h) Pentium 4 in October 2004. This feature is required by all versions of Windows 8. Early Intel 64 implementations had a 36-bit (64 GiB) physical addressing of memory while original AMD64 implementations had a 40-bit (1 TiB) physical addressing. Intel used the 40-bit physical addressing first on Xeon MP (Potomac), launched on 29 March 2005. The difference is not a difference of the user-visible ISAs. In 2007 AMD 10h-based Opteron was the first to provide a 48-bit (256 TiB) physical address space. Intel 64's physical addressing was extended to 44 bits (16 TiB) in Nehalem-EX in 2010 and to 46 bits (64 TiB) in Sandy Bridge E in 2011. With the Ice Lake 3rd gen Xeon Scalable processors, Intel increased the virtual addressing to 57 bits (128 PiB) and physical to 52 bits (4 PiB) in 2021, necessitating a 5-level paging. The following year AMD64 added the same in 4th generation EPYC (Genoa). Non-server CPUs retain smaller address spaces for longer. On all AMD64 processors, the BSF and BSR instructions will, when given a source value of 0, leave their destination register unmodified. This is mostly the case on Intel 64 processors as well, except that on some older Intel 64 CPUs, executing these instructions with an operand size of 32 bits will clear the top 32 bits of their destination register even with a source value of 0 (with the low 32 bits kept unchanged.) AMD64 processors since Opteron Rev. E and Athlon 64 Rev. D reintroduced limited support for segmentation, via the Long Mode Segment Limit Enable (LMSLE) bit, to ease virtualization of 64-bit guests. LMLSE support was removed in the Zen 3 processor. On all Intel 64 processors, CLFLUSH is ordered with respect to SFENCE - this is also the case on newer AMD64 processors (Zen 1 and later). On older AMD64 processors, imposing ordering on the CLFLUSH instruction instead required MFENCE. Adoption In supercomputers tracked by TOP500, the appearance of 64-bit extensions for the x86 architecture enabled 64-bit x86 processors by AMD and Intel to replace most RISC processor architectures previously used in such systems (including PA-RISC, SPARC, Alpha and others), as well as 32-bit x86, even though Intel itself initially tried unsuccessfully to replace x86 with a new incompatible 64-bit architecture in the Itanium processor. , a HPE EPYC-based supercomputer called Frontier is number one. The first ARM-based supercomputer appeared on the list in 2018 and, in recent years, non-CPU architecture co-processors (GPGPU) have also played a big role in performance. Intel's Xeon Phi "Knights Corner" coprocessors, which implement a subset of x86-64 with some vector extensions, are also used, along with x86-64 processors, in the Tianhe-2 supercomputer. Operating system compatibility and characteristics The following operating systems and releases support the x86-64 architecture in long mode. BSD DragonFly BSD Preliminary infrastructure work was started in February 2004 for a x86-64 port. This development later stalled. Development started again during July 2007 and continued during Google Summer of Code 2008 and SoC 2009. The first official release to contain x86-64 support was version 2.4. FreeBSD FreeBSD first added x86-64 support under the name "amd64" as an experimental architecture in 5.1-RELEASE in June 2003. It was included as a standard distribution architecture as of 5.2-RELEASE in January 2004. Since then, FreeBSD has designated it as a Tier 1 platform. The 6.0-RELEASE version cleaned up some quirks with running x86 executables under amd64, and most drivers work just as they do on the x86 architecture. Work is currently being done to integrate more fully the x86 application binary interface (ABI), in the same manner as the Linux 32-bit ABI compatibility currently works. NetBSD x86-64 architecture support was first committed to the NetBSD source tree on June 19, 2001. As of NetBSD 2.0, released on December 9, 2004, NetBSD/amd64 is a fully integrated and supported port. 32-bit code is still supported in 64-bit mode, with a netbsd-32 kernel compatibility layer for 32-bit syscalls. The NX bit is used to provide non-executable stack and heap with per-page granularity (segment granularity being used on 32-bit x86). OpenBSD OpenBSD has supported AMD64 since OpenBSD 3.5, released on May 1, 2004. Complete in-tree implementation of AMD64 support was achieved prior to the hardware's initial release because AMD had loaned several machines for the project's hackathon that year. OpenBSD developers have taken to the platform because of its support for the NX bit, which allowed for an easy implementation of the W^X feature. The code for the AMD64 port of OpenBSD also runs on Intel 64 processors which contains cloned use of the AMD64 extensions, but since Intel left out the page table NX bit in early Intel 64 processors, there is no W^X capability on those Intel CPUs; later Intel 64 processors added the NX bit under the name "XD bit". Symmetric multiprocessing (SMP) works on OpenBSD's AMD64 port, starting with release 3.6 on November 1, 2004. DOS It is possible to enter long mode under DOS without a DOS extender, but the user must return to real mode in order to call BIOS or DOS interrupts. It may also be possible to enter long mode with a DOS extender similar to DOS/4GW, but more complex since x86-64 lacks virtual 8086 mode. DOS itself is not aware of that, and no benefits should be expected unless running DOS in an emulation with an adequate virtualization driver backend, for example: the mass storage interface. Linux Linux was the first operating system kernel to run the x86-64 architecture in long mode, starting with the 2.4 version in 2001 (preceding the hardware's availability). Linux also provides backward compatibility for running 32-bit executables. This permits programs to be recompiled into long mode while retaining the use of 32-bit programs. Current Linux distributions ship with x86-64-native kernels and userlands. Some, such as Arch Linux, SUSE, Mandriva, and Debian, allow users to install a set of 32-bit components and libraries when installing off a 64-bit distribution medium, thus allowing most existing 32-bit applications to run alongside the 64-bit OS. x32 ABI (Application Binary Interface), introduced in Linux 3.4, allows programs compiled for the x32 ABI to run in the 64-bit mode of x86-64 while only using 32-bit pointers and data fields. Though this limits the program to a virtual address space of 4 GiB it also decreases the memory footprint of the program and in some cases can allow it to run faster. 64-bit Linux allows up to 128 TiB of virtual address space for individual processes, and can address approximately 64 TiB of physical memory, subject to processor and system limitations, or up to 128 PiB (virtual) and 4 PiB (physical) with 5-level paging enabled. macOS Mac OS X 10.4.7 and higher versions of Mac OS X 10.4 run 64-bit command-line tools using the POSIX and math libraries on 64-bit Intel-based machines, just as all versions of Mac OS X 10.4 and 10.5 run them on 64-bit PowerPC machines. No other libraries or frameworks work with 64-bit applications in Mac OS X 10.4. The kernel, and all kernel extensions, are 32-bit only. Mac OS X 10.5 supports 64-bit GUI applications using Cocoa, Quartz, OpenGL, and X11 on 64-bit Intel-based machines, as well as on 64-bit PowerPC machines. All non-GUI libraries and frameworks also support 64-bit applications on those platforms. The kernel, and all kernel extensions, are 32-bit only. Mac OS X 10.6 is the first version of macOS that supports a 64-bit kernel. However, not all 64-bit computers can run the 64-bit kernel, and not all 64-bit computers that can run the 64-bit kernel will do so by default. The 64-bit kernel, like the 32-bit kernel, supports 32-bit applications; both kernels also support 64-bit applications. 32-bit applications have a virtual address space limit of 4 GiB under either kernel. The 64-bit kernel does not support 32-bit kernel extensions, and the 32-bit kernel does not support 64-bit kernel extensions. OS X 10.8 includes only the 64-bit kernel, but continues to support 32-bit applications; it does not support 32-bit kernel extensions, however. macOS 10.15 includes only the 64-bit kernel and no longer supports 32-bit applications. This removal of support has presented a problem for WineHQ (and the commercial version CrossOver), as it needs to still be able to run 32-bit Windows applications. The solution, termed wine32on64, was to add thunks that bring the CPU in and out of 32-bit compatibility mode in the nominally 64-bit application. macOS uses the universal binary format to package 32- and 64-bit versions of application and library code into a single file; the most appropriate version is automatically selected at load time. In Mac OS X 10.6, the universal binary format is also used for the kernel and for those kernel extensions that support both 32-bit and 64-bit kernels. Solaris Solaris 10 and later releases support the x86-64 architecture. For Solaris 10, just as with the SPARC architecture, there is only one operating system image, which contains a 32-bit kernel and a 64-bit kernel; this is labeled as the "x64/x86" DVD-ROM image. The default behavior is to boot a 64-bit kernel, allowing both 64-bit and existing or new 32-bit executables to be run. A 32-bit kernel can also be manually selected, in which case only 32-bit executables will run. The isainfo command can be used to determine if a system is running a 64-bit kernel. For Solaris 11, only the 64-bit kernel is provided. However, the 64-bit kernel supports both 32- and 64-bit executables, libraries, and system calls. Windows x64 editions of Microsoft Windows client and server—Windows XP Professional x64 Edition and Windows Server 2003 x64 Edition—were released in March 2005. Internally they are actually the same build (5.2.3790.1830 SP1), as they share the same source base and operating system binaries, so even system updates are released in unified packages, much in the manner as Windows 2000 Professional and Server editions for x86. Windows Vista, which also has many different editions, was released in January 2007. Windows 7 was released in July 2009. Windows Server 2008 R2 was sold in only x64 and Itanium editions; later versions of Windows Server only offer an x64 edition. Versions of Windows for x64 prior to Windows 8.1 and Windows Server 2012 R2 offer the following: 8 TiB of virtual address space per process, accessible from both user mode and kernel mode, referred to as the user mode address space. An x64 program can use all of this, subject to backing store limits on the system, and provided it is linked with the "large address aware" option, which is present by default. This is a 4096-fold increase over the default 2 GiB user-mode virtual address space offered by 32-bit Windows. 8 TiB of kernel mode virtual address space for the operating system. As with the user mode address space, this is a 4096-fold increase over 32-bit Windows versions. The increased space primarily benefits the file system cache and kernel mode "heaps" (non-paged pool and paged pool). Windows only uses a total of 16 TiB out of the 256 TiB implemented by the processors because early AMD64 processors lacked a CMPXCHG16B instruction. Under Windows 8.1 and Windows Server 2012 R2, both user mode and kernel mode virtual address spaces have been extended to 128 TiB. These versions of Windows will not install on processors that lack the CMPXCHG16B instruction. The following additional characteristics apply to all x64 versions of Windows: Ability to run existing 32-bit applications (.exe programs) and dynamic link libraries (.dlls) using WoW64 if WoW64 is supported on that version. Furthermore, a 32-bit program, if it was linked with the "large address aware" option, can use up to 4 GiB of virtual address space in 64-bit Windows, instead of the default 2 GiB (optional 3 GiB with /3GB boot option and "large address aware" link option) offered by 32-bit Windows. Unlike the use of the /3GB boot option on x86, this does not reduce the kernel mode virtual address space available to the operating system. 32-bit applications can, therefore, benefit from running on x64 Windows even if they are not recompiled for x86-64. Both 32- and 64-bit applications, if not linked with "large address aware", are limited to 2 GiB of virtual address space. Ability to use up to 128 GiB (Windows XP/Vista), 192 GiB (Windows 7), 512 GiB (Windows 8), 1 TiB (Windows Server 2003), 2 TiB (Windows Server 2008/Windows 10), 4 TiB (Windows Server 2012), or 24 TiB (Windows Server 2016/2019) of physical random access memory (RAM). LLP64 data model: in C/C++, "int" and "long" types are 32 bits wide, "long long" is 64 bits, while pointers and types derived from pointers are 64 bits wide. Kernel mode device drivers must be 64-bit versions; there is no way to run 32-bit kernel mode executables within the 64-bit operating system. User mode device drivers can be either 32-bit or 64-bit. 16-bit Windows (Win16) and DOS applications will not run on x86-64 versions of Windows due to the removal of the virtual DOS machine subsystem (NTVDM) which relied upon the ability to use virtual 8086 mode. Virtual 8086 mode cannot be entered while running in long mode. Full implementation of the NX (No Execute) page protection feature. This is also implemented on recent 32-bit versions of Windows when they are started in PAE mode. Instead of FS segment descriptor on x86 versions of the Windows NT family, GS segment descriptor is used to point to two operating system defined structures: Thread Information Block (NT_TIB) in user mode and Processor Control Region (KPCR) in kernel mode. Thus, for example, in user mode GS:0 is the address of the first member of the Thread Information Block. Maintaining this convention made the x86-64 port easier, but required AMD to retain the function of the FS and GS segments in long mode – even though segmented addressing per se is not really used by any modern operating system. Early reports claimed that the operating system scheduler would not save and restore the x87 FPU machine state across thread context switches. Observed behavior shows that this is not the case: the x87 state is saved and restored, except for kernel mode-only threads (a limitation that exists in the 32-bit version as well). The most recent documentation available from Microsoft states that the x87/MMX/3DNow! instructions may be used in long mode, but that they are deprecated and may cause compatibility problems in the future. (3DNow! is no longer available on AMD processors, with the exception of the PREFETCH and PREFETCHW instructions, which are also supported on Intel processors as of Broadwell.) Some components like Jet Database Engine and Data Access Objects will not be ported to 64-bit architectures such as x86-64 and IA-64. Microsoft Visual Studio can compile native applications to target either the x86-64 architecture, which can run only on 64-bit Microsoft Windows, or the IA-32 architecture, which can run as a 32-bit application on 32-bit Microsoft Windows or 64-bit Microsoft Windows in WoW64 emulation mode. Managed applications can be compiled either in IA-32, x86-64 or AnyCPU modes. Software created in the first two modes behave like their IA-32 or x86-64 native code counterparts respectively; When using the AnyCPU mode, however, applications in 32-bit versions of Microsoft Windows run as 32-bit applications, while they run as a 64-bit application in 64-bit editions of Microsoft Windows. Video game consoles The PlayStation 4 and Xbox One use AMD x86-64 processors based on the Jaguar microarchitecture. Firmware and games are written in x86-64 code; no legacy x86 code is involved. The PlayStation 5 and Xbox Series X/S use AMD x86-64 processors based on the Zen 2 microarchitecture. The Steam Deck uses a custom AMD x86-64 accelerated processing unit (APU) based on the Zen 2 microarchitecture. Industry naming conventions Since AMD64 and Intel 64 are substantially similar, many software and hardware products use one vendor-neutral term to indicate their compatibility with both implementations. AMD's original designation for this processor architecture, "x86-64", is still used for this purpose, as is the variant "x86_64". Other companies, such as Microsoft and Sun Microsystems/Oracle Corporation, use the contraction "x64" in marketing material. The term IA-64 refers to the Itanium processor, and should not be confused with x86-64, as it is a completely different instruction set. Many operating systems and products, especially those that introduced x86-64 support prior to Intel's entry into the market, use the term "AMD64" or "amd64" to refer to both AMD64 and Intel 64. amd64 Most BSD systems such as FreeBSD, MidnightBSD, NetBSD and OpenBSD refer to both AMD64 and Intel 64 under the architecture name "amd64". Some Linux distributions such as Debian, Ubuntu, Gentoo Linux refer to both AMD64 and Intel 64 under the architecture name "amd64". Microsoft Windows's x64 versions use the AMD64 moniker internally to designate various components which use or are compatible with this architecture. For example, the environment variable PROCESSOR_ARCHITECTURE is assigned the value "AMD64" as opposed to "x86" in 32-bit versions, and the system directory on a Windows x64 Edition installation CD-ROM is named "AMD64", in contrast to "i386" in 32-bit versions. Sun's Solaris's isalist command identifies both AMD64- and Intel 64-based systems as "amd64". Java Development Kit (JDK): the name "amd64" is used in directory names containing x86-64 files. x86_64 The Linux kernel and the GNU Compiler Collection refers to 64-bit architecture as "x86_64". Some Linux distributions, such as Fedora, openSUSE, Arch Linux, Gentoo Linux refer to this 64-bit architecture as "x86_64". Apple macOS refers to 64-bit architecture as "x86-64" or "x86_64", as seen in the Terminal command arch and in their developer documentation. Breaking with most other BSD systems, DragonFly BSD refers to 64-bit architecture as "x86_64". Haiku refers to 64-bit architecture as "x86_64". Licensing x86-64/AMD64 was solely developed by AMD. Until April 2021 when the relevant patents expired, AMD held patents on techniques used in AMD64; those patents had to be licensed from AMD in order to implement AMD64. Intel entered into a cross-licensing agreement with AMD, licensing to AMD their patents on existing x86 techniques, and licensing from AMD their patents on techniques used in x86-64. In 2009, AMD and Intel settled several lawsuits and cross-licensing disagreements, extending their cross-licensing agreements.	Technology	Computer architecture concepts	null
244391	https://en.wikipedia.org/wiki/Flash%20flood	Flash flood	A flash flood is a rapid flooding of low-lying areas: washes, rivers, dry lakes and depressions. It may be caused by heavy rain associated with a severe thunderstorm, hurricane, or tropical storm, or by meltwater from ice and snow. Flash floods may also occur after the collapse of a natural ice or debris dam, or a human structure such as a man-made dam, as occurred before the Johnstown Flood of 1889. Flash floods are distinguished from regular floods by having a timescale of fewer than six hours between rainfall and the onset of flooding. Flash floods are a significant hazard, causing more fatalities in the U.S. in an average year than lightning, tornadoes, or hurricanes. They can also deposit large quantities of sediments on floodplains and destroy vegetation cover not adapted to frequent flood conditions. Causes Flash floods most often occur in dry areas that have recently received precipitation, but they may be seen anywhere downstream from the source of the precipitation, even many miles from the source. In areas on or near volcanoes, flash floods have also occurred after eruptions, when glaciers have been melted by the intense heat. Flash floods are known to occur in the highest mountain ranges of the United States and are also common in the arid plains of the Southwestern United States. Flash flooding can also be caused by extensive rainfall released by hurricanes and other tropical storms, as well as the sudden thawing effect of ice dams. Human activities can also cause flash floods to occur. When dams fail, a large quantity of water can be released and destroy everything in its path. Hazards The United States National Weather Service gives the advice "Turn Around, Don't Drown" for flash floods; that is, it recommends that people get out of the area of a flash flood, rather than trying to cross it. Many people tend to underestimate the dangers of flash floods. What makes flash floods most dangerous is their sudden nature and fast-moving water. A vehicle provides little to no protection against being swept away; it may make people overconfident and less likely to avoid the flash flood. More than half of the fatalities attributed to flash floods are people swept away in vehicles when trying to cross flooded intersections. As little as of water is enough to carry away most SUV-sized vehicles. The U.S. National Weather Service reported in 2005 that, using a national 30-year average, more people die yearly in floods, 127 on average, than by lightning (73), tornadoes (65), or hurricanes (16). In deserts, flash floods can be particularly deadly for several reasons. First, storms in arid regions are infrequent, but they can deliver an enormous amount of water in a very short time. Second, these rains often fall on poorly absorbent and often clay-like soil, which greatly increases the amount of runoff that rivers and other water channels have to handle. These regions tend not to have the infrastructure that wetter regions have to divert water from structures and roads, such as storm drains, culverts, and retention basins, either because of sparse population or poverty, or because residents believe the risk of flash floods is not high enough to justify the expense. In fact, in some areas, desert roads frequently cross a dry river and creek beds without bridges. From the driver's perspective, there may be clear weather, when a river unexpectedly forms ahead of or around the vehicle in a matter of seconds. Finally, the lack of regular rain to clear water channels may cause flash floods in deserts to be headed by large amounts of debris, such as rocks, branches, and logs. Deep slot canyons can be especially dangerous to hikers as they may be flooded by a storm that occurs on a mesa miles away. The flood sweeps through the canyon; the canyon makes it difficult to climb up and out of the way to avoid the flood. For example, a cloudburst in southern Utah on 14 September 2015 resulted in 20 flash flood fatalities, of which seven fatalities occurred at Zion National Park when hikers were trapped by floodwaters in a slot canyon. Impacts Flash floods induce severe impacts in both the built and the natural environment. The effects of flash floods can be catastrophic and show extensive diversity, ranging from damages in buildings and infrastructure to impacts on vegetation, human lives and livestock. The effects are particularly difficult to characterize in urban areas. Researchers have used datasets such as the Severe Hazards Analysis and Verification Experiment (SHAVE) and the U.S. National Weather Service (NWS) Storm Data datasets to connect the impact of flash floods with the physical processes involved in flash flooding. This should increase the reliability of flash flood impact forecasting models. Analysis of flash floods in the United States between 2006 and 2012 shows that injuries and fatalities are most likely in small, rural catchments, that the shortest events are also the most dangerous, that the hazards are greatest after nightfall, and that a very high fraction of injuries and fatalities involve vehicles. An impact severity scale is proposed in 2020 providing a coherent overview of the flash flood effects through the classification of impact types and severity and mapping their spatial extent in a continuous way across the floodplain. Depending on the affected elements, the flood effects are grouped into 4 categories: (i) impacts on built environment (ii) impacts on man-made mobile objects,(iii) impacts on the natural environment (including vegetation, agriculture, geomorphology, and pollution) and (iv) impacts on the human population (entrapments, injuries, fatalities). The scale was proposed as a tool on prevention planning, as the resulting maps offer insights on future impacts, highlighting the high severity areas. Flash floods can cause rapid soil erosion. Much of the Nile delta sedimentation may come from flash flooding in the desert areas that drain into the Nile River. However, flash floods of short duration produce relatively little bedrock erosion or channel widening, having their greatest impact from sedimentation on the floodplain. Some wetlands plants, such as certain varieties of rice, are adapted to endure flash flooding. However, plants that thrive in drier areas can be harmed by flooding, as the plants can become stressed by the large amount of water.	Physical sciences	Water: General	Earth science
244431	https://en.wikipedia.org/wiki/Bipartite%20graph	Bipartite graph	In the mathematical field of graph theory, a bipartite graph (or bigraph) is a graph whose vertices can be divided into two disjoint and independent sets and , that is, every edge connects a vertex in to one in . Vertex sets and are usually called the parts of the graph. Equivalently, a bipartite graph is a graph that does not contain any odd-length cycles. The two sets and may be thought of as a coloring of the graph with two colors: if one colors all nodes in blue, and all nodes in red, each edge has endpoints of differing colors, as is required in the graph coloring problem. In contrast, such a coloring is impossible in the case of a non-bipartite graph, such as a triangle: after one node is colored blue and another red, the third vertex of the triangle is connected to vertices of both colors, preventing it from being assigned either color. One often writes to denote a bipartite graph whose partition has the parts and , with denoting the edges of the graph. If a bipartite graph is not connected, it may have more than one bipartition; in this case, the notation is helpful in specifying one particular bipartition that may be of importance in an application. If , that is, if the two subsets have equal cardinality, then is called a balanced bipartite graph. If all vertices on the same side of the bipartition have the same degree, then is called biregular. Examples When modelling relations between two different classes of objects, bipartite graphs very often arise naturally. For instance, a graph of football players and clubs, with an edge between a player and a club if the player has played for that club, is a natural example of an affiliation network, a type of bipartite graph used in social network analysis. Another example where bipartite graphs appear naturally is in the (NP-complete) railway optimization problem, in which the input is a schedule of trains and their stops, and the goal is to find a set of train stations as small as possible such that every train visits at least one of the chosen stations. This problem can be modeled as a dominating set problem in a bipartite graph that has a vertex for each train and each station and an edge for each pair of a station and a train that stops at that station. A third example is in the academic field of numismatics. Ancient coins are made using two positive impressions of the design (the obverse and reverse). The charts numismatists produce to represent the production of coins are bipartite graphs. More abstract examples include the following: Every tree is bipartite. Cycle graphs with an even number of vertices are bipartite. Every planar graph whose faces all have even length is bipartite. Special cases of this are grid graphs and squaregraphs, in which every inner face consists of 4 edges and every inner vertex has four or more neighbors. The complete bipartite graph on m and n vertices, denoted by Kn,m is the bipartite graph , where U and V are disjoint sets of size m and n, respectively, and E connects every vertex in U with all vertices in V. It follows that Km,n has mn edges. Closely related to the complete bipartite graphs are the crown graphs, formed from complete bipartite graphs by removing the edges of a perfect matching. Hypercube graphs, partial cubes, and median graphs are bipartite. In these graphs, the vertices may be labeled by bitvectors, in such a way that two vertices are adjacent if and only if the corresponding bitvectors differ in a single position. A bipartition may be formed by separating the vertices whose bitvectors have an even number of ones from the vertices with an odd number of ones. Trees and squaregraphs form examples of median graphs, and every median graph is a partial cube. Properties Characterization Bipartite graphs may be characterized in several different ways: An undirected graph is bipartite if and only if it does not contain an odd cycle. A graph is bipartite if and only if it is 2-colorable, (i.e. its chromatic number is less than or equal to 2). A graph is bipartite if and only if every edge belongs to an odd number of bonds, minimal subsets of edges whose removal increases the number of components of the graph. A graph is bipartite if and only if the spectrum of the graph is symmetric. Kőnig's theorem and perfect graphs In bipartite graphs, the size of minimum vertex cover is equal to the size of the maximum matching; this is Kőnig's theorem. An alternative and equivalent form of this theorem is that the size of the maximum independent set plus the size of the maximum matching is equal to the number of vertices. In any graph without isolated vertices the size of the minimum edge cover plus the size of a maximum matching equals the number of vertices. Combining this equality with Kőnig's theorem leads to the facts that, in bipartite graphs, the size of the minimum edge cover is equal to the size of the maximum independent set, and the size of the minimum edge cover plus the size of the minimum vertex cover is equal to the number of vertices. Another class of related results concerns perfect graphs: every bipartite graph, the complement of every bipartite graph, the line graph of every bipartite graph, and the complement of the line graph of every bipartite graph, are all perfect. Perfection of bipartite graphs is easy to see (their chromatic number is two and their maximum clique size is also two) but perfection of the complements of bipartite graphs is less trivial, and is another restatement of Kőnig's theorem. This was one of the results that motivated the initial definition of perfect graphs. Perfection of the complements of line graphs of perfect graphs is yet another restatement of Kőnig's theorem, and perfection of the line graphs themselves is a restatement of an earlier theorem of Kőnig, that every bipartite graph has an edge coloring using a number of colors equal to its maximum degree. According to the strong perfect graph theorem, the perfect graphs have a forbidden graph characterization resembling that of bipartite graphs: a graph is bipartite if and only if it has no odd cycle as a subgraph, and a graph is perfect if and only if it has no odd cycle or its complement as an induced subgraph. The bipartite graphs, line graphs of bipartite graphs, and their complements form four out of the five basic classes of perfect graphs used in the proof of the strong perfect graph theorem. It follows that any subgraph of a bipartite graph is also bipartite because it cannot gain an odd cycle. Degree For a vertex, the number of adjacent vertices is called the degree of the vertex and is denoted . The degree sum formula for a bipartite graph states that The degree sequence of a bipartite graph is the pair of lists each containing the degrees of the two parts and . For example, the complete bipartite graph K3,5 has degree sequence . Isomorphic bipartite graphs have the same degree sequence. However, the degree sequence does not, in general, uniquely identify a bipartite graph; in some cases, non-isomorphic bipartite graphs may have the same degree sequence. The bipartite realization problem is the problem of finding a simple bipartite graph with the degree sequence being two given lists of natural numbers. (Trailing zeros may be ignored since they are trivially realized by adding an appropriate number of isolated vertices to the digraph.) Relation to hypergraphs and directed graphs The biadjacency matrix of a bipartite graph is a (0,1) matrix of size that has a one for each pair of adjacent vertices and a zero for nonadjacent vertices. Biadjacency matrices may be used to describe equivalences between bipartite graphs, hypergraphs, and directed graphs. A hypergraph is a combinatorial structure that, like an undirected graph, has vertices and edges, but in which the edges may be arbitrary sets of vertices rather than having to have exactly two endpoints. A bipartite graph may be used to model a hypergraph in which is the set of vertices of the hypergraph, is the set of hyperedges, and contains an edge from a hypergraph vertex to a hypergraph edge exactly when is one of the endpoints of . Under this correspondence, the biadjacency matrices of bipartite graphs are exactly the incidence matrices of the corresponding hypergraphs. As a special case of this correspondence between bipartite graphs and hypergraphs, any multigraph (a graph in which there may be two or more edges between the same two vertices) may be interpreted as a hypergraph in which some hyperedges have equal sets of endpoints, and represented by a bipartite graph that does not have multiple adjacencies and in which the vertices on one side of the bipartition all have degree two. A similar reinterpretation of adjacency matrices may be used to show a one-to-one correspondence between directed graphs (on a given number of labeled vertices, allowing self-loops) and balanced bipartite graphs, with the same number of vertices on both sides of the bipartition. For, the adjacency matrix of a directed graph with vertices can be any (0,1) matrix of size , which can then be reinterpreted as the adjacency matrix of a bipartite graph with vertices on each side of its bipartition. In this construction, the bipartite graph is the bipartite double cover of the directed graph. Algorithms Testing bipartiteness It is possible to test whether a graph is bipartite, and to return either a two-coloring (if it is bipartite) or an odd cycle (if it is not) in linear time, using depth-first search. The main idea is to assign to each vertex the color that differs from the color of its parent in the depth-first search forest, assigning colors in a preorder traversal of the depth-first-search forest. This will necessarily provide a two-coloring of the spanning forest consisting of the edges connecting vertices to their parents, but it may not properly color some of the non-forest edges. In a depth-first search forest, one of the two endpoints of every non-forest edge is an ancestor of the other endpoint, and when the depth first search discovers an edge of this type it should check that these two vertices have different colors. If they do not, then the path in the forest from ancestor to descendant, together with the miscolored edge, form an odd cycle, which is returned from the algorithm together with the result that the graph is not bipartite. However, if the algorithm terminates without detecting an odd cycle of this type, then every edge must be properly colored, and the algorithm returns the coloring together with the result that the graph is bipartite. Alternatively, a similar procedure may be used with breadth-first search in place of depth-first search. Again, each node is given the opposite color to its parent in the search forest, in breadth-first order. If, when a vertex is colored, there exists an edge connecting it to a previously-colored vertex with the same color, then this edge together with the paths in the breadth-first search forest connecting its two endpoints to their lowest common ancestor forms an odd cycle. If the algorithm terminates without finding an odd cycle in this way, then it must have found a proper coloring, and can safely conclude that the graph is bipartite. For the intersection graphs of line segments or other simple shapes in the Euclidean plane, it is possible to test whether the graph is bipartite and return either a two-coloring or an odd cycle in time , even though the graph itself may have up to edges. Odd cycle transversal Odd cycle transversal is an NP-complete algorithmic problem that asks, given a graph G = (V,E) and a number k, whether there exists a set of k vertices whose removal from G would cause the resulting graph to be bipartite. The problem is fixed-parameter tractable, meaning that there is an algorithm whose running time can be bounded by a polynomial function of the size of the graph multiplied by a larger function of k. The name odd cycle transversal comes from the fact that a graph is bipartite if and only if it has no odd cycles. Hence, to delete vertices from a graph in order to obtain a bipartite graph, one needs to "hit all odd cycle", or find a so-called odd cycle transversal set. In the illustration, every odd cycle in the graph contains the blue (the bottommost) vertices, so removing those vertices kills all odd cycles and leaves a bipartite graph. The edge bipartization problem is the algorithmic problem of deleting as few edges as possible to make a graph bipartite and is also an important problem in graph modification algorithmics. This problem is also fixed-parameter tractable, and can be solved in time , where k is the number of edges to delete and m is the number of edges in the input graph. Matching A matching in a graph is a subset of its edges, no two of which share an endpoint. Polynomial time algorithms are known for many algorithmic problems on matchings, including maximum matching (finding a matching that uses as many edges as possible), maximum weight matching, and stable marriage. In many cases, matching problems are simpler to solve on bipartite graphs than on non-bipartite graphs, and many matching algorithms such as the Hopcroft–Karp algorithm for maximum cardinality matching work correctly only on bipartite inputs. As a simple example, suppose that a set of people are all seeking jobs from among a set of jobs, with not all people suitable for all jobs. This situation can be modeled as a bipartite graph where an edge connects each job-seeker with each suitable job. A perfect matching describes a way of simultaneously satisfying all job-seekers and filling all jobs; Hall's marriage theorem provides a characterization of the bipartite graphs which allow perfect matchings. The National Resident Matching Program applies graph matching methods to solve this problem for U.S. medical student job-seekers and hospital residency jobs. The Dulmage–Mendelsohn decomposition is a structural decomposition of bipartite graphs that is useful in finding maximum matchings. Additional applications Bipartite graphs are extensively used in modern coding theory, especially to decode codewords received from the channel. Factor graphs and Tanner graphs are examples of this. A Tanner graph is a bipartite graph in which the vertices on one side of the bipartition represent digits of a codeword, and the vertices on the other side represent combinations of digits that are expected to sum to zero in a codeword without errors. A factor graph is a closely related belief network used for probabilistic decoding of LDPC and turbo codes. In computer science, a Petri net is a mathematical modeling tool used in analysis and simulations of concurrent systems. A system is modeled as a bipartite directed graph with two sets of nodes: A set of "place" nodes that contain resources, and a set of "event" nodes which generate and/or consume resources. There are additional constraints on the nodes and edges that constrain the behavior of the system. Petri nets utilize the properties of bipartite directed graphs and other properties to allow mathematical proofs of the behavior of systems while also allowing easy implementation of simulations of the system. In projective geometry, Levi graphs are a form of bipartite graph used to model the incidences between points and lines in a configuration. Corresponding to the geometric property of points and lines that every two lines meet in at most one point and every two points be connected with a single line, Levi graphs necessarily do not contain any cycles of length four, so their girth must be six or more.	Mathematics	Graph theory	null
244437	https://en.wikipedia.org/wiki/Hamiltonian%20path	Hamiltonian path	In the mathematical field of graph theory, a Hamiltonian path (or traceable path) is a path in an undirected or directed graph that visits each vertex exactly once. A Hamiltonian cycle (or Hamiltonian circuit) is a cycle that visits each vertex exactly once. A Hamiltonian path that starts and ends at adjacent vertices can be completed by adding one more edge to form a Hamiltonian cycle, and removing any edge from a Hamiltonian cycle produces a Hamiltonian path. The computational problems of determining whether such paths and cycles exist in graphs are NP-complete; see Hamiltonian path problem for details. Hamiltonian paths and cycles are named after William Rowan Hamilton, who invented the icosian game, now also known as Hamilton's puzzle, which involves finding a Hamiltonian cycle in the edge graph of the dodecahedron. Hamilton solved this problem using the icosian calculus, an algebraic structure based on roots of unity with many similarities to the quaternions (also invented by Hamilton). This solution does not generalize to arbitrary graphs. Despite being named after Hamilton, Hamiltonian cycles in polyhedra had also been studied a year earlier by Thomas Kirkman, who, in particular, gave an example of a polyhedron without Hamiltonian cycles. Even earlier, Hamiltonian cycles and paths in the knight's graph of the chessboard, the knight's tour, had been studied in the 9th century in Indian mathematics by Rudrata, and around the same time in Islamic mathematics by al-Adli ar-Rumi. In 18th century Europe, knight's tours were published by Abraham de Moivre and Leonhard Euler. Definitions A Hamiltonian path or traceable path is a path that visits each vertex of the graph exactly once. A graph that contains a Hamiltonian path is called a traceable graph. A graph is Hamiltonian-connected if for every pair of vertices there is a Hamiltonian path between the two vertices. A Hamiltonian cycle, Hamiltonian circuit, vertex tour or graph cycle is a cycle that visits each vertex exactly once. A graph that contains a Hamiltonian cycle is called a Hamiltonian graph. Similar notions may be defined for directed graphs, where each edge (arc) of a path or cycle can only be traced in a single direction (i.e., the vertices are connected with arrows and the edges traced "tail-to-head"). A Hamiltonian decomposition is an edge decomposition of a graph into Hamiltonian circuits. A Hamilton maze is a type of logic puzzle in which the goal is to find the unique Hamiltonian cycle in a given graph. Examples A complete graph with more than two vertices is Hamiltonian Every cycle graph is Hamiltonian Every tournament has an odd number of Hamiltonian paths (Rédei 1934) Every platonic solid, considered as a graph, is Hamiltonian The Cayley graph of a finite Coxeter group is Hamiltonian (For more information on Hamiltonian paths in Cayley graphs, see the Lovász conjecture.) Cayley graphs on nilpotent groups with cyclic commutator subgroup are Hamiltonian. The flip graph of a convex polygon or equivalently, the rotation graph of binary trees, is Hamiltonian. Properties Any Hamiltonian cycle can be converted to a Hamiltonian path by removing one of its edges, but a Hamiltonian path can be extended to a Hamiltonian cycle only if its endpoints are adjacent. All Hamiltonian graphs are biconnected, but a biconnected graph need not be Hamiltonian (see, for example, the Petersen graph). An Eulerian graph (a connected graph in which every vertex has even degree) necessarily has an Euler tour, a closed walk passing through each edge of exactly once. This tour corresponds to a Hamiltonian cycle in the line graph , so the line graph of every Eulerian graph is Hamiltonian. Line graphs may have other Hamiltonian cycles that do not correspond to Euler tours, and in particular the line graph of every Hamiltonian graph is itself Hamiltonian, regardless of whether the graph is Eulerian. A tournament (with more than two vertices) is Hamiltonian if and only if it is strongly connected. The number of different Hamiltonian cycles in a complete undirected graph on vertices is and in a complete directed graph on vertices is . These counts assume that cycles that are the same apart from their starting point are not counted separately. Bondy–Chvátal theorem The best vertex degree characterization of Hamiltonian graphs was provided in 1972 by the Bondy–Chvátal theorem, which generalizes earlier results by G. A. Dirac (1952) and Øystein Ore. Both Dirac's and Ore's theorems can also be derived from Pósa's theorem (1962). Hamiltonicity has been widely studied with relation to various parameters such as graph density, toughness, forbidden subgraphs and distance among other parameters. Dirac and Ore's theorems basically state that a graph is Hamiltonian if it has enough edges. The Bondy–Chvátal theorem operates on the closure of a graph with vertices, obtained by repeatedly adding a new edge connecting a nonadjacent pair of vertices and with until no more pairs with this property can be found. As complete graphs are Hamiltonian, all graphs whose closure is complete are Hamiltonian, which is the content of the following earlier theorems by Dirac and Ore. The following theorems can be regarded as directed versions: The number of vertices must be doubled because each undirected edge corresponds to two directed arcs and thus the degree of a vertex in the directed graph is twice the degree in the undirected graph. The above theorem can only recognize the existence of a Hamiltonian path in a graph and not a Hamiltonian Cycle. Many of these results have analogues for balanced bipartite graphs, in which the vertex degrees are compared to the number of vertices on a single side of the bipartition rather than the number of vertices in the whole graph. Existence of Hamiltonian cycles in planar graphs The Hamiltonian cycle polynomial An algebraic representation of the Hamiltonian cycles of a given weighted digraph (whose arcs are assigned weights from a certain ground field) is the Hamiltonian cycle polynomial of its weighted adjacency matrix defined as the sum of the products of the arc weights of the digraph's Hamiltonian cycles. This polynomial is not identically zero as a function in the arc weights if and only if the digraph is Hamiltonian. The relationship between the computational complexities of computing it and computing the permanent was shown by Grigoriy Kogan.	Mathematics	Graph theory	null
244439	https://en.wikipedia.org/wiki/Turnstone	Turnstone	Turnstones are two bird species that constitute the genus Arenaria in the family Scolopacidae. They are closely related to calidrid sandpipers and might be considered members of the tribe Calidriini. The genus Arenaria was introduced by the French zoologist Mathurin Jacques Brisson in 1760 with the ruddy turnstone (Arenaria interpres) as the type species. The genus name arenaria is from Latin arenarius, "inhabiting sand", from arena, "sand". The genus contains two species: the ruddy turnstone (Arenaria interpres) and the black turnstone (Arenaria melanocephala). Both birds are waders. Their length is typically between 20 and 25 cm, with a wingspan between 50 and 60 cm and a body mass between 110 and 130g. For waders their build is stocky, with short, slightly upturned, wedge shaped bills. They have white patches on the back, wings and tail. They are high Arctic breeders, and are migratory. Their strong necks and powerful, slightly upturned bills are adapted to their feeding technique. As the name implies, these species overturn stones, seaweed, and similar items in search of invertebrate prey. They are strictly coastal, prefer stony beaches to sand, and often share beach space with other species of waders such as purple sandpipers. Species There exists a fossil bone, a distal piece of tarsometatarsus found in the Edson Beds of Sherman County, Kansas. Dating from the mid-Blancan some 4-3 million years ago, it appears to be from a calidriid somewhat similar to a pectoral sandpiper, but has some traits reminiscent of turnstones. Depending on which traits are apomorphic and plesiomorphic, it may be an ancestral representative of either lineage.	Biology and health sciences	Charadriiformes	Animals
244463	https://en.wikipedia.org/wiki/Adjacency%20matrix	Adjacency matrix	In graph theory and computer science, an adjacency matrix is a square matrix used to represent a finite graph. The elements of the matrix indicate whether pairs of vertices are adjacent or not in the graph. In the special case of a finite simple graph, the adjacency matrix is a (0,1)-matrix with zeros on its diagonal. If the graph is undirected (i.e. all of its edges are bidirectional), the adjacency matrix is symmetric. The relationship between a graph and the eigenvalues and eigenvectors of its adjacency matrix is studied in spectral graph theory. The adjacency matrix of a graph should be distinguished from its incidence matrix, a different matrix representation whose elements indicate whether vertex–edge pairs are incident or not, and its degree matrix, which contains information about the degree of each vertex. Definition For a simple graph with vertex set , the adjacency matrix is a square matrix such that its element is 1 when there is an edge from vertex to vertex , and 0 when there is no edge. The diagonal elements of the matrix are all 0, since edges from a vertex to itself (loops) are not allowed in simple graphs. It is also sometimes useful in algebraic graph theory to replace the nonzero elements with algebraic variables. The same concept can be extended to multigraphs and graphs with loops by storing the number of edges between each two vertices in the corresponding matrix element, and by allowing nonzero diagonal elements. Loops may be counted either once (as a single edge) or twice (as two vertex-edge incidences), as long as a consistent convention is followed. Undirected graphs often use the latter convention of counting loops twice, whereas directed graphs typically use the former convention. Of a bipartite graph The adjacency matrix of a bipartite graph whose two parts have and vertices can be written in the form where is an matrix, and and represent the and zero matrices. In this case, the smaller matrix uniquely represents the graph, and the remaining parts of can be discarded as redundant. is sometimes called the biadjacency matrix. Formally, let be a bipartite graph with parts , and edges . The biadjacency matrix is the 0–1 matrix in which if and only if . If is a bipartite multigraph or weighted graph, then the elements are taken to be the number of edges between the vertices or the weight of the edge , respectively. Variations An -adjacency matrix of a simple graph has if is an edge, if it is not, and on the diagonal. The Seidel adjacency matrix is a -adjacency matrix. This matrix is used in studying strongly regular graphs and two-graphs. The distance matrix has in position the distance between vertices and . The distance is the length of a shortest path connecting the vertices. Unless lengths of edges are explicitly provided, the length of a path is the number of edges in it. The distance matrix resembles a high power of the adjacency matrix, but instead of telling only whether or not two vertices are connected (i.e., the connection matrix, which contains Boolean values), it gives the exact distance between them. Examples Undirected graphs The convention followed here (for undirected graphs) is that each edge adds 1 to the appropriate cell in the matrix, and each loop (an edge from a vertex to itself) adds 2 to the appropriate cell on the diagonal in the matrix. This allows the degree of a vertex to be easily found by taking the sum of the values in either its respective row or column in the adjacency matrix. Directed graphs The adjacency matrix of a directed graph can be asymmetric. One can define the adjacency matrix of a directed graph either such that a non-zero element indicates an edge from to or it indicates an edge from to . The former definition is commonly used in graph theory and social network analysis (e.g., sociology, political science, economics, psychology). The latter is more common in other applied sciences (e.g., dynamical systems, physics, network science) where is sometimes used to describe linear dynamics on graphs. Using the first definition, the in-degrees of a vertex can be computed by summing the entries of the corresponding column and the out-degree of vertex by summing the entries of the corresponding row. When using the second definition, the in-degree of a vertex is given by the corresponding row sum and the out-degree is given by the corresponding column sum. Trivial graphs The adjacency matrix of a complete graph contains all ones except along the diagonal where there are only zeros. The adjacency matrix of an empty graph is a zero matrix. Properties Spectrum The adjacency matrix of an undirected simple graph is symmetric, and therefore has a complete set of real eigenvalues and an orthogonal eigenvector basis. The set of eigenvalues of a graph is the spectrum of the graph. It is common to denote the eigenvalues by The greatest eigenvalue is bounded above by the maximum degree. This can be seen as result of the Perron–Frobenius theorem, but it can be proved easily. Let be one eigenvector associated to and the entry in which has maximum absolute value. Without loss of generality assume is positive since otherwise you simply take the eigenvector -, also associated to . Then For -regular graphs, is the first eigenvalue of for the vector (it is easy to check that it is an eigenvalue and it is the maximum because of the above bound). The multiplicity of this eigenvalue is the number of connected components of , in particular for connected graphs. It can be shown that for each eigenvalue , its opposite is also an eigenvalue of if is a bipartite graph. In particular − is an eigenvalue of any -regular bipartite graph. The difference is called the spectral gap and it is related to the expansion of . It is also useful to introduce the spectral radius of denoted by . This number is bounded by . This bound is tight in the Ramanujan graphs, which have applications in many areas. Isomorphism and invariants Suppose two directed or undirected graphs and with adjacency matrices and are given. and are isomorphic if and only if there exists a permutation matrix such that In particular, and are similar and therefore have the same minimal polynomial, characteristic polynomial, eigenvalues, determinant and trace. These can therefore serve as isomorphism invariants of graphs. However, two graphs may possess the same set of eigenvalues but not be isomorphic. Such linear operators are said to be isospectral. Matrix powers If is the adjacency matrix of the directed or undirected graph , then the matrix (i.e., the matrix product of copies of ) has an interesting interpretation: the element gives the number of (directed or undirected) walks of length from vertex to vertex . If is the smallest nonnegative integer, such that for some , , the element of is positive, then is the distance between vertex and vertex . A great example of how this is useful is in counting the number of triangles in an undirected graph , which is exactly the trace of divided by 3 or 6 depending on whether the graph is directed or not. We divide by those values to compensate for the overcounting of each triangle. In an undirected graph, each triangle will be counted twice for all three nodes, because the path can be followed clockwise or counterclockwise : ijk or ikj. The adjacency matrix can be used to determine whether or not the graph is connected. If a directed graph has a nilpotent adjacency matrix (i.e., if there exists such that is the zero matrix), then it is a directed acyclic graph. Data structures The adjacency matrix may be used as a data structure for the representation of graphs in computer programs for manipulating graphs. The main alternative data structure, also in use for this application, is the adjacency list. The space needed to represent an adjacency matrix and the time needed to perform operations on them is dependent on the matrix representation chosen for the underlying matrix. Sparse matrix representations only store non-zero matrix entries and implicitly represent the zero entries. They can, for example, be used to represent sparse graphs without incurring the space overhead from storing the many zero entries in the adjacency matrix of the sparse graph. In the following section the adjacency matrix is assumed to be represented by an array data structure so that zero and non-zero entries are all directly represented in storage. Because each entry in the adjacency matrix requires only one bit, it can be represented in a very compact way, occupying only \|V \|2&hairsp;/&hairsp;8 bytes to represent a directed graph, or (by using a packed triangular format and only storing the lower triangular part of the matrix) approximately \|V \|2&hairsp;/&hairsp;16 bytes to represent an undirected graph. Although slightly more succinct representations are possible, this method gets close to the information-theoretic lower bound for the minimum number of bits needed to represent all -vertex graphs. For storing graphs in text files, fewer bits per byte can be used to ensure that all bytes are text characters, for instance by using a Base64 representation. Besides avoiding wasted space, this compactness encourages locality of reference. However, for a large sparse graph, adjacency lists require less storage space, because they do not waste any space representing edges that are not present. An alternative form of adjacency matrix (which, however, requires a larger amount of space) replaces the numbers in each element of the matrix with pointers to edge objects (when edges are present) or null pointers (when there is no edge). It is also possible to store edge weights directly in the elements of an adjacency matrix. Besides the space tradeoff, the different data structures also facilitate different operations. Finding all vertices adjacent to a given vertex in an adjacency list is as simple as reading the list, and takes time proportional to the number of neighbors. With an adjacency matrix, an entire row must instead be scanned, which takes a larger amount of time, proportional to the number of vertices in the whole graph. On the other hand, testing whether there is an edge between two given vertices can be determined at once with an adjacency matrix, while requiring time proportional to the minimum degree of the two vertices with the adjacency list.	Mathematics	Graph theory	null
244517	https://en.wikipedia.org/wiki/Tibia	Tibia	The tibia (; : tibiae or tibias), also known as the shinbone or shankbone, is the larger, stronger, and anterior (frontal) of the two bones in the leg below the knee in vertebrates (the other being the fibula, behind and to the outside of the tibia); it connects the knee with the ankle. The tibia is found on the medial side of the leg next to the fibula and closer to the median plane. The tibia is connected to the fibula by the interosseous membrane of leg, forming a type of fibrous joint called a syndesmosis with very little movement. The tibia is named for the flute tibia. It is the second largest bone in the human body, after the femur. The leg bones are the strongest long bones as they support the rest of the body. Structure In human anatomy, the tibia is the second largest bone next to the femur. As in other vertebrates the tibia is one of two bones in the lower leg, the other being the fibula, and is a component of the knee and ankle joints. The ossification or formation of the bone starts from three centers, one in the shaft and one in each extremity. The tibia is categorized as a long bone and is as such composed of a diaphysis and two epiphyses. The diaphysis is the midsection of the tibia, also known as the shaft or body. While the epiphyses are the two rounded extremities of the bone; an upper (also known as superior or proximal) closest to the thigh and a lower (also known as inferior or distal) closest to the foot. The tibia is most contracted in the lower third and the distal extremity is smaller than the proximal. Upper extremity Condyles of tibia The proximal or upper extremity of the tibia is expanded in the transverse plane with a medial and lateral condyle, which are both flattened in the horizontal plane. The medial condyle is the larger of the two and is better supported over the shaft. The upper surfaces of the condyles articulate with the femur to form the tibiofemoral joint, the weightbearing part of the knee joint. The medial and lateral condyle are separated by the intercondylar area, where the cruciate ligaments and the menisci attach. Here the medial and lateral intercondylar tubercle forms the intercondylar eminence. Together with the medial and lateral condyle the intercondylar region forms the tibial plateau, which both articulates with and is anchored to the lower extremity of the femur. The intercondylar eminence divides the intercondylar area into an anterior and posterior part. The anterolateral region of the anterior intercondylar area are perforated by numerous small openings for nutrient arteries. The articular surfaces of both condyles are concave, particularly centrally. The flatter outer margins are in contact with the menisci. The medial condyles superior surface is oval in form and extends laterally onto the side of medial intercondylar tubercle. The lateral condyles superior surface is more circular in form and its medial edge extends onto the side of the lateral intercondylar tubercle. The posterior surface of the medial condyle bears a horizontal groove for part of the attachment of the semimembranosus muscle, whereas the lateral condyle has a circular facet for articulation with the head of the fibula. Beneath the condyles is the tibial tuberosity which serves for attachment of the patellar ligament, a continuation of the quadriceps femoris muscle. Facets The superior articular surface presents two smooth articular facets. The medial facet, oval in shape, is slightly concave from side to side, and from before backward. The lateral, nearly circular, is concave from side to side, but slightly convex from before backward, especially at its posterior part, where it is prolonged on to the posterior surface for a short distance. The central portions of these facets articulate with the condyles of the femur, while their peripheral portions support the menisci of the knee joint, which here intervene between the two bones. Intercondyloid eminence Between the articular facets in the intercondylar area, but nearer the posterior than the anterior aspect of the bone, is the intercondyloid eminence (spine of tibia), surmounted on either side by a prominent tubercle, on to the sides of which the articular facets are prolonged; in front of and behind the intercondyloid eminence are rough depressions for the attachment of the anterior and posterior cruciate ligaments and the menisci. Surfaces The anterior surfaces of the condyles are continuous with one another, forming a large somewhat flattened area; this area is triangular, broad above, and perforated by large vascular foramina; narrow below where it ends in a large oblong elevation, the tuberosity of the tibia, which gives attachment to the patellar ligament; a bursa intervenes between the deep surface of the ligament and the part of the bone immediately above the tuberosity. Posteriorly, the condyles are separated from each other by a shallow depression, the posterior intercondyloid fossa, which gives attachment to part of the posterior cruciate ligament of the knee-joint. The medial condyle presents posteriorly a deep transverse groove, for the insertion of the tendon of the semimembranosus. Its medial surface is convex, rough, and prominent; it gives attachment to the medial collateral ligament. The lateral condyle presents posteriorly a flat articular facet, nearly circular in form, directed downward, backward, and lateralward, for articulation with the head of the fibula. Its lateral surface is convex, rough, and prominent in front: on it is an eminence, situated on a level with the upper border of the tuberosity and at the junction of its anterior and lateral surfaces, for the attachment of the iliotibial band. Just below this a part of the extensor digitorum longus takes origin and a slip from the tendon of the biceps femoris is inserted. Shaft The shaft or body of the tibia is triangular in cross-section and forms three borders: an anterior, medial, and lateral or interosseous border. These three borders form three surfaces: the medial, lateral, and posterior. Borders The anterior crest or border, the most prominent of the three, commences above at the tuberosity, and ends below at the anterior margin of the medial malleolus. It is sinuous and prominent in the upper two-thirds of its extent, but smooth and rounded below; it gives attachment to the deep fascia of the leg. The medial border is smooth and rounded above and below, but more prominent in the center. It begins at the back part of the medial condyle, and ends at the posterior border of the medial malleolus; its upper part gives attachment to the tibial collateral ligament of the knee-joint to the extent of about 5 cm., and insertion to some fibers of the popliteus muscle. From its middle third some fibers of the soleus and flexor digitorum longus muscles take origin. The interosseous crest or lateral border is thin and prominent, especially its central part, and gives attachment to the interosseous membrane; it commences above in front of the fibular articular facet, and bifurcates below, to form the boundaries of a triangular rough surface, for the attachment of the interosseous ligament connecting the tibia and fibula. Surfaces The medial surface is smooth, convex, and broader above than below; its upper third, directed forward and medialward, is covered by the aponeurosis derived from the tendon of the sartorius, and by the tendons of the Gracilis and Semitendinosus, all of which are inserted nearly as far forward as the anterior crest; in the rest of its extent it is subcutaneous. The lateral surface is narrower than the medial; its upper two-thirds present a shallow groove for the origin of the Tibialis anterior; its lower third is smooth, convex, curves gradually forward to the anterior aspect of the bone, and is covered by the tendons of the Tibialis anterior, Extensor hallucis longus, and Extensor digitorum longus, arranged in this order from the medial side. The posterior surface presents, at its upper part, a prominent ridge, the popliteal line, which extends obliquely downward from the back part of the articular facet for the fibula to the medial border, at the junction of its upper and middle thirds; it marks the lower limit of the insertion of the Popliteus, serves for the attachment of the fascia covering this muscle, and gives origin to part of the Soleus, Flexor digitorum longus, and Tibialis posterior. The triangular area, above this line, gives insertion to the Popliteus. The middle third of the posterior surface is divided by a vertical ridge into two parts; the ridge begins at the popliteal line and is well-marked above, but indistinct below; the medial and broader portion gives origin to the Flexor digitorum longus, the lateral and narrower to part of the Tibialis posterior. The remaining part of the posterior surface is smooth and covered by the Tibialis posterior, Flexor digitorum longus, and Flexor hallucis longus. Immediately below the popliteal line is the nutrient foramen, which is large and directed obliquely downward. Lower extremity The distal end of the tibia is much smaller than the proximal end and presents five surfaces; it is prolonged downward on its medial side as a strong pyramidal process, the medial malleolus. The lower extremity of the tibia together with the fibula and talus forms the ankle joint. Surfaces The inferior articular surface is quadrilateral, and smooth for articulation with the talus. It is concave from before backward, broader in front than behind, and traversed from before backward by a slight elevation, separating two depressions. It is continuous with that on the medial malleolus. The anterior surface of the lower extremity is smooth and rounded above, and covered by the tendons of the Extensor muscles; its lower margin presents a rough transverse depression for the attachment of the articular capsule of the ankle-joint. The posterior surface is traversed by a shallow groove directed obliquely downward and medialward, continuous with a similar groove on the posterior surface of the talus and serving for the passage of the tendon of the Flexor hallucis longus. The lateral surface presents a triangular rough depression for the attachment of the inferior interosseous ligament connecting it with the fibula; the lower part of this depression is smooth, covered with cartilage in the fresh state, and articulates with the fibula. The surface is bounded by two prominent borders (the anterior and posterior colliculi), continuous above with the interosseous crest; they afford attachment to the anterior and posterior ligaments of the lateral malleolus. The medial surface – see medial malleolus for details. Fractures Ankle fractures of the tibia have several classification systems based on location or mechanism: Medial malleolus – Herscovici classification Posterior malleolus – Haruguchi classification Mechanism – Lauge-Hansen classification Blood supply The tibia is supplied with blood from two sources: A nutrient artery, as the main source, and periosteal vessels derived from the anterior tibial artery. Joints The tibia is a part of four joints; the knee, ankle, superior and inferior tibiofibular joint. In the knee the tibia forms one of the two articulations with the femur, often referred to as the tibiofemoral components of the knee joint.; it is the weightbearing part of the knee joint. The tibiofibular joints are the articulations between the tibia and fibula which allows very little movement. The proximal tibiofibular joint is a small plane joint. The joint is formed between the undersurface of the lateral tibial condyle and the head of fibula. The joint capsule is reinforced by anterior and posterior ligament of the head of the fibula. The distal tibiofibular joint (tibiofibular syndesmosis) is formed by the rough, convex surface of the medial side of the distal end of the fibula, and a rough concave surface on the lateral side of the tibia. The part of the ankle joint known as the talocrural joint, is a synovial hinge joint that connects the distal ends of the tibia and fibula in the lower limb with the proximal end of the talus. The articulation between the tibia and the talus bears more weight than between the smaller fibula and the talus. Development The tibia is ossified from three centers: a primary center for the diaphysis (shaft) and a secondary center for each epiphysis (extremity). Ossification begins in the center of the body, about the seventh week of fetal life, and gradually extends toward the extremities. The center for the upper epiphysis appears before or shortly after birth at close to 34 weeks gestation; it is flattened in form, and has a thin tongue-shaped process in front, which forms the tuberosity; that for the lower epiphysis appears in the second year. The lower epiphysis fuses with the tibial shaft at about the eighteenth, and the upper one fuses about the twentieth year. Two additional centers occasionally exist, one for the tongue-shaped process of the upper epiphysis, which forms the tuberosity, and one for the medial malleolus. Function Muscle attachments Strength The tibia has been modeled as taking an axial force during walking that is up to 4.7 bodyweight. Its bending moment in the sagittal plane in the late stance phase is up to 71.6 bodyweight times millimetre. Clinical significance Fracture Fractures of the tibia can be divided into those that only involve the tibia; bumper fracture, Segond fracture, Gosselin fracture, toddler's fracture, and those including both the tibia and fibula; trimalleolar fracture, bimalleolar fracture, Pott's fracture. Society and culture In Judaism, the tibia, or shankbone, of a goat or sheep is used in the Passover Seder plate. Other animals The structure of the tibia in most other tetrapods is essentially similar to that in humans. The tuberosity of the tibia, a crest to which the patellar ligament attaches in mammals, is instead the point for the tendon of the quadriceps muscle in reptiles, birds, and amphibians, which have no patella. Additional images	Biology and health sciences	Skeletal system	Biology
244611	https://en.wikipedia.org/wiki/Newton%27s%20law%20of%20universal%20gravitation	Newton's law of universal gravitation	Newton's law of universal gravitation states that every particle attracts every other particle in the universe with a force that is proportional to the product of their masses and inversely proportional to the square of the distance between their centers. Separated objects attract and are attracted as if all their mass were concentrated at their centers. The publication of the law has become known as the "first great unification", as it marked the unification of the previously described phenomena of gravity on Earth with known astronomical behaviors. This is a general physical law derived from empirical observations by what Isaac Newton called inductive reasoning.<ref>Isaac Newton: "In [experimental] philosophy particular propositions are inferred from the phenomena and afterwards rendered general by induction": Principia', Book 3, General Scholium, at p.392 in Volume 2 of Andrew Motte's English translation published 1729.</ref> It is a part of classical mechanics and was formulated in Newton's work Philosophiæ Naturalis Principia Mathematica ("the Principia"), first published on 5 July 1687. The equation for universal gravitation thus takes the form: where F is the gravitational force acting between two objects, m1 and m2 are the masses of the objects, r is the distance between the centers of their masses, and G is the gravitational constant. The first test of Newton's law of gravitation between masses in the laboratory was the Cavendish experiment conducted by the British scientist Henry Cavendish in 1798. It took place 111 years after the publication of Newton's Principia and approximately 71 years after his death. Newton's law of gravitation resembles Coulomb's law of electrical forces, which is used to calculate the magnitude of the electrical force arising between two charged bodies. Both are inverse-square laws, where force is inversely proportional to the square of the distance between the bodies. Coulomb's law has charge in place of mass and a different constant. Newton's law was later superseded by Albert Einstein's theory of general relativity, but the universality of the gravitational constant is intact and the law still continues to be used as an excellent approximation of the effects of gravity in most applications. Relativity is required only when there is a need for extreme accuracy, or when dealing with very strong gravitational fields, such as those found near extremely massive and dense objects, or at small distances (such as Mercury's orbit around the Sun). History Before Newton’s law of gravity, there were many theories explaining gravity. Philoshophers made observations about things falling down − and developed theories why they do – as early as Aristotle who thought that rocks fall to the ground because seeking the ground was an essential part of their nature. Around 1600, the scientific method began to take root. René Descartes started over with a more fundamental view, developing ideas of matter and action independent of theology. Galileo Galilei wrote about experimental measurements of falling and rolling objects. Johannes Kepler's laws of planetary motion summarized Tycho Brahe's astronomical observations. Around 1666 Isaac Newton developed the idea that Kepler's laws must also apply to the orbit of the Moon around the Earth and then to all objects on Earth. The analysis required assuming that the gravitation force acted as if all of the mass of the Earth were concentrated at its center, an unproven conjecture at that time. His calculations of the Moon orbit time was within 16% of the known value. By 1680, new values for the diameter of the Earth improved his orbit time to within 1.6%, but more importantly Newton had found a proof of his earlier conjecture. In 1687 Newton published his Principia which combined his laws of motion with new mathematical analysis to explain Kepler's empirical results. His explanation was in the form of a law of universal gravitation: any two bodies are attracted by a force proportional to their mass and inversely proportional to their separation squared. Newton's original formula was: where the symbol means "is proportional to". To make this into an equal-sided formula or equation, there needed to be a multiplying factor or constant that would give the correct force of gravity no matter the value of the masses or distance between them (the gravitational constant). Newton would need an accurate measure of this constant to prove his inverse-square law. When Newton presented Book 1 of the unpublished text in April 1686 to the Royal Society, Robert Hooke made a claim that Newton had obtained the inverse square law from him, ultimately a frivolous accusation. Newton's "causes hitherto unknown" While Newton was able to formulate his law of gravity in his monumental work, he was deeply uncomfortable with the notion of "action at a distance" that his equations implied. In 1692, in his third letter to Bentley, he wrote: "That one body may act upon another at a distance through a vacuum without the mediation of anything else, by and through which their action and force may be conveyed from one another, is to me so great an absurdity that, I believe, no man who has in philosophic matters a competent faculty of thinking could ever fall into it." He never, in his words, "assigned the cause of this power". In all other cases, he used the phenomenon of motion to explain the origin of various forces acting on bodies, but in the case of gravity, he was unable to experimentally identify the motion that produces the force of gravity (although he invented two mechanical hypotheses in 1675 and 1717). Moreover, he refused to even offer a hypothesis as to the cause of this force on grounds that to do so was contrary to sound science. He lamented that "philosophers have hitherto attempted the search of nature in vain" for the source of the gravitational force, as he was convinced "by many reasons" that there were "causes hitherto unknown" that were fundamental to all the "phenomena of nature". These fundamental phenomena are still under investigation and, though hypotheses abound, the definitive answer has yet to be found. And in Newton's 1713 General Scholium in the second edition of Principia: "I have not yet been able to discover the cause of these properties of gravity from phenomena and I feign no hypotheses. ... It is enough that gravity does really exist and acts according to the laws I have explained, and that it abundantly serves to account for all the motions of celestial bodies." Modern form In modern language, the law states the following: Assuming SI units, F is measured in newtons (N), m1 and m2 in kilograms (kg), r in meters (m), and the constant G is The value of the constant G was first accurately determined from the results of the Cavendish experiment conducted by the British scientist Henry Cavendish in 1798, although Cavendish did not himself calculate a numerical value for G. This experiment was also the first test of Newton's theory of gravitation between masses in the laboratory. It took place 111 years after the publication of Newton's Principia and 71 years after Newton's death, so none of Newton's calculations could use the value of G; instead he could only calculate a force relative to another force. Bodies with spatial extent If the bodies in question have spatial extent (as opposed to being point masses), then the gravitational force between them is calculated by summing the contributions of the notional point masses that constitute the bodies. In the limit, as the component point masses become "infinitely small", this entails integrating the force (in vector form, see below) over the extents of the two bodies. In this way, it can be shown that an object with a spherically symmetric distribution of mass exerts the same gravitational attraction on external bodies as if all the object's mass were concentrated at a point at its center. (This is not generally true for non-spherically symmetrical bodies.) For points inside a spherically symmetric distribution of matter, Newton's shell theorem can be used to find the gravitational force. The theorem tells us how different parts of the mass distribution affect the gravitational force measured at a point located a distance r0 from the center of the mass distribution: The portion of the mass that is located at radii causes the same force at the radius r0 as if all of the mass enclosed within a sphere of radius r0 was concentrated at the center of the mass distribution (as noted above). The portion of the mass that is located at radii exerts no net gravitational force at the radius r0 from the center. That is, the individual gravitational forces exerted on a point at radius r0 by the elements of the mass outside the radius r0 cancel each other. As a consequence, for example, within a shell of uniform thickness and density there is no net gravitational acceleration anywhere within the hollow sphere. Vector form Newton's law of universal gravitation can be written as a vector equation to account for the direction of the gravitational force as well as its magnitude. In this formula, quantities in bold represent vectors. where F21 is the force applied on body 2 exerted by body 1, G is the gravitational constant, m1 and m2 are respectively the masses of bodies 1 and 2, r21 = r2 − r1 is the displacement vector between bodies 1 and 2, and is the unit vector from body 1 to body 2. It can be seen that the vector form of the equation is the same as the scalar form given earlier, except that F is now a vector quantity, and the right hand side is multiplied by the appropriate unit vector. Also, it can be seen that F12 = −F21. Gravity field The gravitational field is a vector field that describes the gravitational force that would be applied on an object in any given point in space, per unit mass. It is actually equal to the gravitational acceleration at that point. It is a generalisation of the vector form, which becomes particularly useful if more than two objects are involved (such as a rocket between the Earth and the Moon). For two objects (e.g. object 2 is a rocket, object 1 the Earth), we simply write r instead of r12 and m instead of m2 and define the gravitational field g(r) as: so that we can write: This formulation is dependent on the objects causing the field. The field has units of acceleration; in SI, this is m/s2. Gravitational fields are also conservative; that is, the work done by gravity from one position to another is path-independent. This has the consequence that there exists a gravitational potential field V(r) such that If m1 is a point mass or the mass of a sphere with homogeneous mass distribution, the force field g(r) outside the sphere is isotropic, i.e., depends only on the distance r from the center of the sphere. In that case As per Gauss's law, field in a symmetric body can be found by the mathematical equation: where is a closed surface and is the mass enclosed by the surface. Hence, for a hollow sphere of radius and total mass , For a uniform solid sphere of radius and total mass , Limitations Newton's description of gravity is sufficiently accurate for many practical purposes and is therefore widely used. Deviations from it are small when the dimensionless quantities and are both much less than one, where is the gravitational potential, is the velocity of the objects being studied, and is the speed of light in vacuum. For example, Newtonian gravity provides an accurate description of the Earth/Sun system, since where is the radius of the Earth's orbit around the Sun. In situations where either dimensionless parameter is large, then general relativity must be used to describe the system. General relativity reduces to Newtonian gravity in the limit of small potential and low velocities, so Newton's law of gravitation is often said to be the low-gravity limit of general relativity. Observations conflicting with Newton's formula Newton's theory does not fully explain the precession of the perihelion of the orbits of the planets, especially that of Mercury, which was detected long after the life of Newton. There is a 43 arcsecond per century discrepancy between the Newtonian calculation, which arises only from the gravitational attractions from the other planets, and the observed precession, made with advanced telescopes during the 19th century. The predicted angular deflection of light rays by gravity (treated as particles travelling at the expected speed) that is calculated by using Newton's theory is only one-half of the deflection that is observed by astronomers. Calculations using general relativity are in much closer agreement with the astronomical observations. In spiral galaxies, the orbiting of stars around their centers seems to strongly disobey both Newton's law of universal gravitation and general relativity. Astrophysicists, however, explain this marked phenomenon by assuming the presence of large amounts of dark matter. Einstein's solution The first two conflicts with observations above were explained by Einstein's theory of general relativity, in which gravitation is a manifestation of curved spacetime instead of being due to a force propagated between bodies. In Einstein's theory, energy and momentum distort spacetime in their vicinity, and other particles move in trajectories determined by the geometry of spacetime. This allowed a description of the motions of light and mass that was consistent with all available observations. In general relativity, the gravitational force is a fictitious force resulting from the curvature of spacetime, because the gravitational acceleration of a body in free fall is due to its world line being a geodesic of spacetime. Extensions In recent years, quests for non-inverse square terms in the law of gravity have been carried out by neutron interferometry. Solutions The two-body problem has been completely solved, as has the restricted three-body problem. The n-body problem is an ancient, classical problem of predicting the individual motions of a group of celestial objects interacting with each other gravitationally. Solving this problem – from the time of the Greeks and on – has been motivated by the desire to understand the motions of the Sun, planets and the visible stars. The classical problem can be informally stated as: given the quasi-steady orbital properties (instantaneous position, velocity and time) of a group of celestial bodies, predict their interactive forces; and consequently, predict their true orbital motions for all future times. In the 20th century, understanding the dynamics of globular cluster star systems became an important n-body problem too. The n''-body problem in general relativity is considerably more difficult to solve.	Physical sciences	Classical mechanics	null

Subsets and Splits

No community queries yet

The top public SQL queries from the community will appear here once available.