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455672 | https://en.wikipedia.org/wiki/Color%20theory | Color theory | Color theory, or more specifically traditional color theory, is a historical body of knowledge describing the behavior of colors, namely in color mixing, color contrast effects, color harmony, color schemes and color symbolism. Modern color theory is generally referred to as color science. While there is no clear distinction in scope, traditional color theory tends to be more subjective and have artistic applications, while color science tends to be more objective and have functional applications, such as in chemistry, astronomy or color reproduction. Color theory dates back at least as far as Aristotle's treatise On Colors. A formalization of "color theory" began in the 18th century, initially within a partisan controversy over Isaac Newton's theory of color (Opticks, 1704) and the nature of primary colors. By the end of the 19th century, a schism had formed between traditional color theory and color science.
History
Color theory is rooted in antiquity, with early musings on color in Aristotle's (d. 322 BCE) On Colors and Claudius Ptolemy's (d. 168 CE) Optics. The influence of light on color was investigated and revealed further by al-Kindi (d. 873) and Ibn al-Haytham (d. 1039). Ibn Sina (d. 1037), Nasir al-Din al-Tusi (d. 1274), and Robert Grosseteste (d. 1253) discovered that contrary to the teachings of Aristotle, there are multiple color paths to get from black to white. More modern approaches to color theory principles can be found in the writings of Leone Battista Alberti (c. 1435) and the notebooks of Leonardo da Vinci (c. 1490).
The RYB primary colors became the foundation of 18th-century theories of color vision, as the fundamental sensory qualities that are blended in the perception of all physical colors, and conversely, in the physical mixture of pigments or dyes. These theories were enhanced by 18th-century investigations of a variety of purely psychological color effects, in particular the contrast between "complementary" or opposing hues that are produced by color afterimages and in the contrasting shadows in colored light. These ideas and many personal color observations were summarized in two founding documents in color theory: the Theory of Colours (1810) by the German poet Johann Wolfgang von Goethe, and The Law of Simultaneous Color Contrast (1839) by the French industrial chemist Michel Eugène Chevreul. Charles Hayter published A New Practical Treatise on the Three Primitive Colours Assumed as a Perfect System of Rudimentary Information (London 1826), in which he described how all colors could be obtained from just three.
Subsequently, German and English scientists established in the late 19th century that color perception is best described in terms of a different set of primary colors—red, green and blue-violet (RGB)—modeled through the additive mixture of three monochromatic lights. Subsequent research anchored these primary colors in the differing responses to light by three types of color receptors or cones in the retina (trichromacy). On this basis the quantitative description of the color mixture or colorimetry developed in the early 20th century, along with a series of increasingly sophisticated models of color space and color perception, such as the opponent process theory.
Across the same period, industrial chemistry radically expanded the color range of lightfast synthetic pigments, allowing for substantially improved saturation in color mixtures of dyes, paints, and inks. It also created the dyes and chemical processes necessary for color photography. As a result, three-color printing became aesthetically and economically feasible in mass printed media, and the artists' color theory was adapted to primary colors most effective in inks or photographic dyes: cyan, magenta, and yellow (CMY). (In printing, dark colors are supplemented by black ink, called "key," to make the CMYK system; in both printing and photography, white is provided by the color of the paper.) These CMY primary colors were reconciled with the RGB primaries, and subtractive color mixing with additive color mixing, by defining the CMY primaries as substances that absorbed only one of the retinal primary colors: cyan absorbs only red (−R+G+B), magenta only green (+R−G+B), and yellow only blue-violet (+R+G−B). It is important to add that the CMYK, or process, color printing is meant as an economical way of producing a wide range of colors for printing, but is deficient in reproducing certain colors, notably orange and slightly deficient in reproducing purples. A wider range of colors can be obtained with the addition of other colors to the printing process, such as in Pantone's Hexachrome printing ink system (six colors), among others.
For much of the 19th century artistic color theory either lagged behind scientific understanding or was augmented by science books written for the lay public, in particular Modern Chromatics (1879) by the American physicist Ogden Rood, and early color atlases developed by Albert Munsell (Munsell Book of Color, 1915, see Munsell color system) and Wilhelm Ostwald (Color Atlas, 1919). Major advances were made in the early 20th century by artists teaching or associated with the German Bauhaus, in particular Wassily Kandinsky, Johannes Itten, Faber Birren and Josef Albers, whose writings mix speculation with an empirical or demonstration-based study of color design principles.
Color mixing
One of the earliest purposes of color theory was to establish rules governing the mixing of pigments.
Traditional color theory was built around "pure" or ideal colors, characterized by different sensory experiences rather than attributes of the physical world. This has led to several inaccuracies in traditional color theory principles that are not always remedied in modern formulations. Another issue has been the tendency to describe color effects holistically or categorically, for example as a contrast between "yellow" and "blue" conceived as generic colors instead of the three color attributes generally considered by color science: hue, colorfulness and lightness. These confusions are partly historical and arose in scientific uncertainty about color perception that was not resolved until the late 19th century when artistic notions were already entrenched. They also arise from the attempt to describe the highly contextual and flexible behavior of color perception in terms of abstract color sensations that can be generated equivalently by any visual media.
Primary colors
Color theory asserts three pure primary colors that can be used to mix all possible colors. These are sometimes considered as red, yellow and blue (RYB) or as red, green and blue (RGB). Ostensibly, any failure of specific paints or inks to match this ideal performance is due to the impurity or imperfection of the colorants. In contrast, modern color science does not recognize universal primary colors (no finite combination of colors can produce all other colors) and only uses primary colors to define a given color space. Any three primary colors can mix only a limited range of colors, called a gamut, which is always smaller (contains fewer colors) than the full range of colors humans can perceive. Primary colors also can't be made from other colors as they are inherently pure and distinct.
Complementary colors
For the mixing of colored light, Isaac Newton's color wheel is often used to describe complementary colors, which are colors that cancel each other's hue to produce an achromatic (white, gray or black) light mixture. Newton offered as a conjecture that colors exactly opposite one another on the hue circle cancel out each other's hue; this concept was demonstrated more thoroughly in the 19th century. An example of complementary colors would be magenta and green.
A key assumption in Newton's hue circle was that the "fiery" or maximum saturated hues are located on the outer circumference of the circle, while achromatic white is at the center. Then the saturation of the mixture of two spectral hues was predicted by the straight line between them; the mixture of three colors was predicted by the "center of gravity" or centroid of three triangle points, and so on.
According to traditional color theory based on subtractive primary colors and the RYB color model, yellow mixed with purple, orange mixed with blue, or red mixed with green produces an equivalent gray and are the painter's complementary colors.
One reason the artist's primary colors work at all is due to the imperfect pigments being used have sloped absorption curves and change color with concentration. A pigment that is pure red at high concentrations can behave more like magenta at low concentrations. This allows it to make purples that would otherwise be impossible. Likewise, a blue that is ultramarine at high concentrations appears cyan at low concentrations, allowing it to be used to mix green. Chromium red pigments can appear orange, and then yellow, as the concentration is reduced. It is even possible to mix very low concentrations of the blue mentioned and the chromium red to get a greenish color. This works much better with oil colors than it does with watercolors and dyes.
The old primaries depend on sloped absorption curves and pigment leakages to work, while newer scientifically derived ones depend solely on controlling the amount of absorption in certain parts of the spectrum.
Tints and shades
When mixing pigments, a color is produced which is always darker and lower in chroma, or saturation, than the parent colors. This moves the mixed color toward a neutral color—a gray or near-black. Lights are made brighter or dimmer by adjusting their brightness, or energy level; in painting, lightness is adjusted through mixture with white, black, or a color's complement.
It is common among some painters to darken a paint color by adding black paint—producing colors called shades—or lighten a color by adding white—producing colors called tints. However, it is not always the best way for representational painting, as an unfortunate result is for colors to also shift in hue. For instance, darkening a color by adding black can cause colors such as yellows, reds, and oranges, to shift toward the greenish or bluish part of the spectrum. Lightening a color by adding white can cause a shift towards blue when mixed with reds and oranges. Another practice when darkening a color is to use its opposite, or complementary, color (e.g. purplish-red added to yellowish-green) to neutralize it without a shift in hue and darken it if the additive color is darker than the parent color. When lightening a color this hue shift can be corrected with the addition of a small amount of an adjacent color to bring the hue of the mixture back in line with the parent color (e.g. adding a small amount of orange to a mixture of red and white will correct the tendency of this mixture to shift slightly towards the blue end of the spectrum).
Split primary palette
The split-primary palette is a color-wheel model that relies on misconceptions to attempt to explain the unsatisfactory results produced when mixing the traditional primary colors, red, yellow, and blue.
Painters have long considered red, yellow, and blue to be primary colors. In practice, however, some of the mixtures produced from these colors lack chromatic intensity. Rather than adopt a more effective set of primary colors, proponents of split-primary theory explain this lack of chroma by the purported presence of impurities, small amounts of other colors in the paints, or biases away from the ideal primary toward one or the other of the adjacent colors. Every red paint, for example, is said to be tainted with, or biased toward, either blue or yellow, every blue paint toward either red or green, and every yellow toward either green or orange. These biases are said to result in mixtures that contain sets of complementary colors, darkening the resulting color. To obtain vivid mixed colors, according to split-primary theory, it is necessary to employ two primary colors whose biases both fall in the direction, on the color wheel, of the color to be mixed, combining, for example, green-biased blue and green-biased yellow to make bright green. Based on this reasoning, proponents of split-primary theory conclude that two versions of each primary color, often called "cool" and "warm," are needed in order to mix a wide gamut of high-chroma colors.
In fact, the perceived bias of colors is not due to impurity. Rather, the appearance of any given colorant is inherent to its chemical and physical properties, and its purity unrelated to whether it conforms to our arbitrary conception of an ideal hue. Moreover, the identity of gamut-optimizing primary colors is determined by the physiology of human color vision. Although no set of three primary paints can be mixed to obtain the complete color gamut perceived by humans, red, yellow, and blue are a poor choice if high-chroma mixtures are desired. This is because painting is a subtractive color process, for which red and blue are secondary, not primary, colors.
Although flawed in principle, the split-primary system can be successful in practice, because the recommended blue-biased red and green-biased blue positions are often filled by near approximations of magenta and cyan, respectively, while orange-biased red and violet-biased blue serve as secondary colors, tending to further widen the mixable gamut.
This system is in effect a simplified version of Newton's geometrical rule that colors closer together on the hue circle will produce more vibrant mixtures. A mixture produced from two primary colors, however, will be much more highly saturated than one produced from two secondary colors, even though the pairs are the same distance apart on the hue circle, revealing the limitations of the circular model in the prediction of color-mixing results. For example, a mixture of magenta and cyan inks or paints will produce vivid blues and violets, whereas a mixture of red and blue inks or paints will produce darkened violets and purples, even though the angular distance separating magenta and cyan is the same as that separating red and blue.
Color contrast
In Chevreul's 1839 book The principles of harmony and contrast of colours, he introduced the law of color contrast, stating that colors that appear together (spatially or temporally) will be altered as if mixed with the complementary color of the other color, functionally boosting the color contrast between them. For example, a piece of yellow fabric placed on a blue background will appear tinted orange because orange is the complementary color to blue. Chevreul formalized three types of contrast:
simultaneous contrast, which appears in two colors viewed side by side,
successive contrast, for the afterimage left on an achromatic background after viewing a color, and
mixed contrast, for the afterimage left on another color.
Warm vs. cool colors
The distinction between "warm" and "cool" colors has been important since at least the late 18th century. The difference (as traced by etymologies in the Oxford English Dictionary), seems related to the observed contrast in landscape light, between the "warm" colors associated with daylight or sunset, and the "cool" colors associated with a gray or overcast day. Warm colors are often said to be hues from red through yellow, browns, and tans included; cool colors are often said to be the hues from blue-green through blue violet, most grays included. There is a historical disagreement about the colors that anchor the polarity, but 19th-century sources put the peak contrast between red-orange and greenish-blue.
Color theory has described perceptual and psychological effects to this contrast. Warm colors are said to advance or appear more active in a painting, while cool colors tend to recede; used in interior design or fashion, warm colors are said to arouse or stimulate the viewer, while cool colors calm and relax. Most of these effects, to the extent they are real, can be attributed to the higher saturation and lighter value of warm pigments in contrast to cool pigments; brown is a dark, unsaturated warm color that few people think of as visually active or psychologically arousing.
Color harmony and color schemes
It has been suggested that "Colors seen together to produce a pleasing affective response are said to be in harmony". However, color harmony is a complex notion because human responses to color are both affective and cognitive, involving emotional response and judgment. Hence, our responses to color and the notion of color harmony is open to the influence of a range of different factors. These factors include individual differences (such as age, gender, personal preference, affective state, etc.) as well as cultural, sub-cultural, and socially-based differences which gives rise to conditioning and learned responses about color. In addition, context always has an influence on responses about color and the notion of color harmony, and this concept is also influenced by temporal factors (such as changing trends) and perceptual factors (such as simultaneous contrast) which may impinge on human response to color. The following conceptual model illustrates this 21st-century approach to color harmony:
wherein color harmony is a function (f) of the interaction between color/s (Col 1, 2, 3, …, n) and the factors that influence positive aesthetic response to color: individual differences (ID) such as age, gender, personality and affective state; cultural experiences (CE), the prevailing context (CX) which includes setting and ambient lighting; intervening perceptual effects (P) and the effects of time (T) in terms of prevailing social trends.
In addition, given that humans can perceive around 2.3 million different colors, it has been suggested that the number of possible color combinations is virtually infinite thereby implying that predictive color harmony formulae are fundamentally unsound. Despite this, many color theorists have devised formulae, principles or guidelines for color combination with the aim being to predict or specify positive aesthetic response or "color harmony".
Color wheel models have often been used as a basis for color combination guidelines and for defining relationships between colors. Some theorists and artists believe juxtapositions of complementary color will produce strong contrast, a sense of visual tension as well as "color harmony"; while others believe juxtapositions of analogous colors will elicit a positive aesthetic response. Color combination guidelines (or formulas) suggest that colors next to each other on the color wheel model (analogous colors) tend to produce a single-hued or monochromatic color experience and some theorists also refer to these as "simple harmonies".
In addition, split complementary color schemes usually depict a modified complementary pair, with instead of the "true" second color being chosen, a range of analogous hues around it are chosen, i.e. the split complements of red are blue-green and yellow-green. A triadic color scheme adopts any three colors approximately equidistant around a color wheel model. Feisner and Mahnke are among a number of authors who provide color combination guidelines in greater detail.
Color combination formulae and principles may provide some guidance but have limited practical application. This is due to the influence of contextual, perceptual, and temporal factors which will influence how color/s are perceived in any given situation, setting, or context. Such formulae and principles may be useful in fashion, interior and graphic design, but much depends on the tastes, lifestyle, and cultural norms of the viewer or consumer.
Black and white have long been known to combine "well" with almost any other colors; black decreases the apparent saturation or brightness of colors paired with it and white shows off all hues to equal effect.
Color symbolism
A major underpinning of traditional color theory is that colors carry significant cultural symbolism, or even have immutable, universal meaning. As early as the ancient Greek philosophers, many theorists have devised color associations and linked particular connotative meanings to specific colors. However, connotative color associations and color symbolism tends to be culture-bound and may also vary across different contexts and circumstances. For example, red has many different connotative and symbolic meanings from exciting, arousing, sensual, romantic, and feminine; to a symbol of good luck; and also acts as a signal of danger. Such color associations tend to be learned and do not necessarily hold irrespective of individual and cultural differences or contextual, temporal or perceptual factors. It is important to note that while color symbolism and color associations exist, their existence does not provide evidential support for color psychology or claims that color has therapeutic properties.
| Physical sciences | Basics_7 | null |
455769 | https://en.wikipedia.org/wiki/Rigid%20body | Rigid body | In physics, a rigid body, also known as a rigid object, is a solid body in which deformation is zero or negligible. The distance between any two given points on a rigid body remains constant in time regardless of external forces or moments exerted on it. A rigid body is usually considered as a continuous distribution of mass.
In the study of special relativity, a perfectly rigid body does not exist; and objects can only be assumed to be rigid if they are not moving near the speed of light. In quantum mechanics, a rigid body is usually thought of as a collection of point masses. For instance, molecules (consisting of the point masses: electrons and nuclei) are often seen as rigid bodies (see classification of molecules as rigid rotors).
Kinematics
Linear and angular position
The position of a rigid body is the position of all the particles of which it is composed. To simplify the description of this position, we exploit the property that the body is rigid, namely that all its particles maintain the same distance relative to each other. If the body is rigid, it is sufficient to describe the position of at least three non-collinear particles. This makes it possible to reconstruct the position of all the other particles, provided that their time-invariant position relative to the three selected particles is known. However, typically a different, mathematically more convenient, but equivalent approach is used. The position of the whole body is represented by:
the linear position or position of the body, namely the position of one of the particles of the body, specifically chosen as a reference point (typically coinciding with the center of mass or centroid of the body), together with
the angular position (also known as orientation, or attitude) of the body.
Thus, the position of a rigid body has two components: linear and angular, respectively. The same is true for other kinematic and kinetic quantities describing the motion of a rigid body, such as linear and angular velocity, acceleration, momentum, impulse, and kinetic energy.
The linear position can be represented by a vector with its tail at an arbitrary reference point in space (the origin of a chosen coordinate system) and its tip at an arbitrary point of interest on the rigid body, typically coinciding with its center of mass or centroid. This reference point may define the origin of a coordinate system fixed to the body.
There are several ways to numerically describe the orientation of a rigid body, including a set of three Euler angles, a quaternion, or a direction cosine matrix (also referred to as a rotation matrix). All these methods actually define the orientation of a basis set (or coordinate system) which has a fixed orientation relative to the body (i.e. rotates together with the body), relative to another basis set (or coordinate system), from which the motion of the rigid body is observed. For instance, a basis set with fixed orientation relative to an airplane can be defined as a set of three orthogonal unit vectors b1, b2, b3, such that b1 is parallel to the chord line of the wing and directed forward, b2 is normal to the plane of symmetry and directed rightward, and b3 is given by the cross product .
In general, when a rigid body moves, both its position and orientation vary with time. In the kinematic sense, these changes are referred to as translation and rotation, respectively. Indeed, the position of a rigid body can be viewed as a hypothetic translation and rotation (roto-translation) of the body starting from a hypothetic reference position (not necessarily coinciding with a position actually taken by the body during its motion).
Linear and angular velocity
Velocity (also called linear velocity) and angular velocity are measured with respect to a frame of reference.
The linear velocity of a rigid body is a vector quantity, equal to the time rate of change of its linear position. Thus, it is the velocity of a reference point fixed to the body. During purely translational motion (motion with no rotation), all points on a rigid body move with the same velocity. However, when motion involves rotation, the instantaneous velocity of any two points on the body will generally not be the same. Two points of a rotating body will have the same instantaneous velocity only if they happen to lie on an axis parallel to the instantaneous axis of rotation.
Angular velocity is a vector quantity that describes the angular speed at which the orientation of the rigid body is changing and the instantaneous axis about which it is rotating (the existence of this instantaneous axis is guaranteed by the Euler's rotation theorem). All points on a rigid body experience the same angular velocity at all times. During purely rotational motion, all points on the body change position except for those lying on the instantaneous axis of rotation. The relationship between orientation and angular velocity is not directly analogous to the relationship between position and velocity. Angular velocity is not the time rate of change of orientation, because there is no such concept as an orientation vector that can be differentiated to obtain the angular velocity.
Kinematical equations
Addition theorem for angular velocity
The angular velocity of a rigid body B in a reference frame N is equal to the sum of the angular velocity of a rigid body D in N and the angular velocity of B with respect to D:
In this case, rigid bodies and reference frames are indistinguishable and completely interchangeable.
Addition theorem for position
For any set of three points P, Q, and R, the position vector from P to R is the sum of the position vector from P to Q and the position vector from Q to R:
The norm of a position vector is the spatial distance.
Here the coordinates of all three vectors must be expressed in coordinate frames with the same orientation.
Mathematical definition of velocity
The velocity of point P in reference frame N is defined as the time derivative in N of the position vector from O to P:
where O is any arbitrary point fixed in reference frame N, and the N to the left of the d/dt operator indicates that the derivative is taken in reference frame N. The result is independent of the selection of O so long as O is fixed in N.
Mathematical definition of acceleration
The acceleration of point P in reference frame N is defined as the time derivative in N of its velocity:
Velocity of two points fixed on a rigid body
For two points P and Q that are fixed on a rigid body B, where B has an angular velocity in the reference frame N, the velocity of Q in N can be expressed as a function of the velocity of P in N:
where is the position vector from P to Q., with coordinates expressed in N (or a frame with the same orientation as N.)
This relation can be derived from the temporal invariance of the norm distance between P and Q.
Acceleration of two points fixed on a rigid body
By differentiating the equation for the Velocity of two points fixed on a rigid body in N with respect to time, the acceleration in reference frame N of a point Q fixed on a rigid body B can be expressed as
where is the angular acceleration of B in the reference frame N.
Angular velocity and acceleration of two points fixed on a rigid body
As mentioned above, all points on a rigid body B have the same angular velocity in a fixed reference frame N, and thus the same angular acceleration
Velocity of one point moving on a rigid body
If the point R is moving in the rigid body B while B moves in reference frame N, then the velocity of R in N is
where Q is the point fixed in B that is instantaneously coincident with R at the instant of interest. This relation is often combined with the relation for the Velocity of two points fixed on a rigid body.
Acceleration of one point moving on a rigid body
The acceleration in reference frame N of the point R moving in body B while B is moving in frame N is given by
where Q is the point fixed in B that instantaneously coincident with R at the instant of interest. This equation is often combined with Acceleration of two points fixed on a rigid body.
Other quantities
If C is the origin of a local coordinate system L, attached to the body,
the spatial or twist acceleration of a rigid body is defined as the spatial acceleration of C (as opposed to material acceleration above):
where
represents the position of the point/particle with respect to the reference point of the body in terms of the local coordinate system L (the rigidity of the body means that this does not depend on time)
is the orientation matrix, an orthogonal matrix with determinant 1, representing the orientation (angular position) of the local coordinate system L, with respect to the arbitrary reference orientation of another coordinate system G. Think of this matrix as three orthogonal unit vectors, one in each column, which define the orientation of the axes of L with respect to G.
represents the angular velocity of the rigid body
represents the total velocity of the point/particle
represents the total acceleration of the point/particle
represents the angular acceleration of the rigid body
represents the spatial acceleration of the point/particle
represents the spatial acceleration of the rigid body (i.e. the spatial acceleration of the origin of L).
In 2D, the angular velocity is a scalar, and matrix A(t) simply represents a rotation in the xy-plane by an angle which is the integral of the angular velocity over time.
Vehicles, walking people, etc., usually rotate according to changes in the direction of the velocity: they move forward with respect to their own orientation. Then, if the body follows a closed orbit in a plane, the angular velocity integrated over a time interval in which the orbit is completed once, is an integer times 360°. This integer is the winding number with respect to the origin of the velocity. Compare the amount of rotation associated with the vertices of a polygon.
Kinetics
Any point that is rigidly connected to the body can be used as reference point (origin of coordinate system L) to describe the linear motion of the body (the linear position, velocity and acceleration vectors depend on the choice).
However, depending on the application, a convenient choice may be:
the center of mass of the whole system, which generally has the simplest motion for a body moving freely in space;
a point such that the translational motion is zero or simplified, e.g. on an axle or hinge, at the center of a ball and socket joint, etc.
When the center of mass is used as reference point:
The (linear) momentum is independent of the rotational motion. At any time it is equal to the total mass of the rigid body times the translational velocity.
The angular momentum with respect to the center of mass is the same as without translation: at any time it is equal to the inertia tensor times the angular velocity. When the angular velocity is expressed with respect to a coordinate system coinciding with the principal axes of the body, each component of the angular momentum is a product of a moment of inertia (a principal value of the inertia tensor) times the corresponding component of the angular velocity; the torque is the inertia tensor times the angular acceleration.
Possible motions in the absence of external forces are translation with constant velocity, steady rotation about a fixed principal axis, and also torque-free precession.
The net external force on the rigid body is always equal to the total mass times the translational acceleration (i.e., Newton's second law holds for the translational motion, even when the net external torque is nonzero, and/or the body rotates).
The total kinetic energy is simply the sum of translational and rotational energy.
Geometry
Two rigid bodies are said to be different (not copies) if there is no proper rotation from one to the other.
A rigid body is called chiral if its mirror image is different in that sense, i.e., if it has either no symmetry or its symmetry group contains only proper rotations. In the opposite case an object is called achiral: the mirror image is a copy, not a different object. Such an object may have a symmetry plane, but not necessarily: there may also be a plane of reflection with respect to which the image of the object is a rotated version. The latter applies for S2n, of which the case n = 1 is inversion symmetry.
For a (rigid) rectangular transparent sheet, inversion symmetry corresponds to having on one side an image without rotational symmetry and on the other side an image such that what shines through is the image at the top side, upside down. We can distinguish two cases:
the sheet surface with the image is not symmetric - in this case the two sides are different, but the mirror image of the object is the same, after a rotation by 180° about the axis perpendicular to the mirror plane.
the sheet surface with the image has a symmetry axis - in this case the two sides are the same, and the mirror image of the object is also the same, again after a rotation by 180° about the axis perpendicular to the mirror plane.
A sheet with a through and through image is achiral. We can distinguish again two cases:
the sheet surface with the image has no symmetry axis - the two sides are different
the sheet surface with the image has a symmetry axis - the two sides are the same
Configuration space
The configuration space of a rigid body with one point fixed (i.e., a body with zero translational motion) is given by the underlying manifold of the rotation group SO(3). The configuration space of a nonfixed (with non-zero translational motion) rigid body is E+(3), the subgroup of direct isometries of the Euclidean group in three dimensions (combinations of translations and rotations).
| Physical sciences | Solid mechanics | Physics |
456161 | https://en.wikipedia.org/wiki/Tropical%20climate | Tropical climate | Tropical climate is the first of the five major climate groups in the Köppen climate classification identified with the letter A. Tropical climates are defined by a monthly average temperature of or higher in the coolest month, featuring hot temperatures and high humidity all year-round. Annual precipitation is often abundant in tropical climates, and shows a seasonal rhythm but may have seasonal dryness to varying degrees. There are normally only two seasons in tropical climates, a wet (rainy/monsoon) season and a dry season. The annual temperature range in tropical climates is normally very small. Sunlight is intense in these climates.
There are three basic types of tropical climates within the tropical climate group: tropical rainforest climate (Af), tropical monsoon climate (Am) and tropical savanna or tropical wet and dry climate (Aw for dry winters, and As for dry summers), which are classified and distinguished by the precipitation levels of the driest month in those regions.
Köppen climate classification
The Köppen climate classification is the most widely used climate classification system. It defines a tropical climate as a region where the mean temperature of the coldest month is greater than or equal to and does not fit into the criteria for B-group climates, classifying them as an A-group (tropical climate group). A-group regions are usually found in the tropics, below 23.5 latitude in both the southern and northern hemisphere; they include areas around the Equator, Central America, North-central portions of South America, central Africa, southern portions of Asia and parts of North Australia and the Pacific Ocean islands.
In Group A, there are three types of this climate: the tropical rainforest climate (Af), tropical monsoon climate (Am) and tropical wet and dry or savanna climate (Aw or As). All of the three climates are classified by their Pdry (short for precipitation of the driest month). Tropical rainforest climate's Pdry should be greater than or equal . Tropical monsoon climate's Pdry should be in the range from to 60 mm. Tropical wet and dry or savanna climate's Pdry should be less than .
Tropical climate biome
Tropical climates normally have only two seasons, a wet season and a dry season. Depending on the location of the region, the wet and dry seasons can have varying duration. Annual temperature changes in the tropics are small. Due to the high temperatures and abundant rainfall, much of the plant life grows throughout the year. High temperature and humidity is the most suitable environment for epiphytes to grow. In many tropical climates, vegetation grow in layers: shrubs under tall trees, bushes under shrubs and grasses under bushes. Tropical plants are rich in resources, including coffee, cocoa and oil palm. Listed below are the types of vegetation unique to each of the three climates that make up the tropical climate biome.
Natural vegetation
Tropical rainforest vegetation including: Bengal bamboo, bougainvillea, curare, coconut tree, durian and banana.
Tropical monsoon vegetation including: teak, deodar, rosewood, sandalwood and bamboo.
Tropical wet and dry or savanna vegetation including: acacia senegal, elephant grass, jarrah tree, gum tree eucalyptus and whistling thorn.
Tropical rainforest climate
The Köppen classification identifies tropical rainforest climates (Zone Af: f = "feucht", German for moist) as usually having north and south latitudinal ranges of just 5-10 degrees from the equator. Tropical rainforest climates have high temperatures: the yearly average temperature is normally between . The precipitation can reach over 100 inches a year. The seasons are evenly distributed throughout the year, and there is almost no drought period here. Regions that contain tropical rainforest climate mainly include the upper Amazon basin of South America, the Northern Zaire (Congo) basin of Africa, and the islands of the East Indies.
The tropical rainforest climate differs from other subtypes of tropical climates as it has more kinds of trees due to its precipitation. The large number of trees contribute back to the humidity of the climate because of the transpiration, which is the process of water evaporated from the surface of living plants to the atmosphere. The warmth and abundant precipitation heavily contributes to the diversity and characteristics of vegetations under the tropical rainforest climate. The vegetations develop a vertical stratification and various growth forms to receive enough sunlight, which is unusual under other types of climate.
Tropical monsoon climate
The Köppen classification tool identifies tropical monsoon climate as having small annual temperature ranges, high temperatures, and plentiful precipitation. This climate also has a short dry season which almost always occurs in the winter. The tropical monsoon climate is often found within countries in the south and southeast Asia region between the latitude of 10 degrees north and the Tropic of Cancer. It can also be found in West Africa and South America. The annual temperature of regions under tropical monsoon climate is also stable.
The tropical monsoon climate has the following main characteristic. The average annual temperature is around and has an average annual temperature range of about . Distinction between wet and drought seasons, the tropical monsoon climate is different from other tropical climates because of its uneven precipitation throughout the year.
There are three main seasons of tropical monsoon climate: the cool dry season is from fall to late winter, the hot dry season is in the spring and the rainy or monsoon season is near or during the summer months.
The tropical monsoon forest mainly consists of three layered structures. The first layer is the surface layer which is a very dense layer of shrubs and grasses. The second layer is the understory layer with trees about 15 meters tall. The top layer is called the canopy tree layer which has trees from 25 to 40 meters tall and those trees grow closely while above is the emergent layer with sporadic trees taller than 35 meters.
Tropical savanna or wet and dry climate
Tropical savanna climates, or tropical wet and dry climates, are mainly located between the 10° and 25° north-south latitudes, and often occur at the outer margins of the tropics. Typical regions include central Africa, parts of South America, as well as northern and eastern Australia. The temperature range of savanna climate is between . In summer, the temperature is between 25 °C and 30 °C, while in winter the temperature is between 20 °C and 30 °C, but still stays above an 18 °C mean. The annual precipitation is between 700 and 1000 mm. The driest months are generally in the winter and they have less than 60 mm of rainfall (often much less).
Regions under the savanna climate usually have lands covered with flat grassland vegetation with areas of woodlands. Those grassland biomes cover almost 20% of the Earth's surface. The grassland vegetation types include Rhodes grass, red oats grass, star grass and lemongrass.
| Physical sciences | Climatology | null |
456169 | https://en.wikipedia.org/wiki/Sodalite | Sodalite | Sodalite ( ) is a tectosilicate mineral with the formula , with royal blue varieties widely used as an ornamental gemstone. Although massive sodalite samples are opaque, crystals are usually transparent to translucent. Sodalite is a member of the sodalite group with hauyne, nosean, lazurite and tugtupite.
The people of the Caral culture traded for sodalite from the Collao altiplano. First discovered by Europeans in 1811 in the Ilimaussaq intrusive complex in Greenland, sodalite did not become widely important as an ornamental stone until 1891 when vast deposits of fine material were discovered in Ontario, Canada.
Structure
The structure of sodalite was first studied by Linus Pauling in 1930. It is a cubic mineral of space group P3n (space group 218) which consists of an aluminosilicate cage network with Na+ cations and chloride anions in the interframework. (There may be small amounts of other cations and anions instead.) This framework forms a zeolite cage structure. Each unit cell has two cavities, which have almost the same structure as the borate cage found in the zinc borate , the beryllosilicate cage , and the aluminate cage in ), and as in the similar mineral tugtupite () (see Haüyne#Sodalite group). There is one cavity around each chloride ion. One chloride is located at the corners of the unit cell, and the other at the centre. Each cavity has chiral tetrahedral symmetry, and the cavities around these two chloride locations are mirror images one of the other (a glide plane or a four-fold improper rotation takes one into the other). There are four sodium ions around each chloride ion (at one distance, and four more at a greater distance), surrounded by twelve tetrahedra and twelve tetrahedra. The silicon and aluminum atoms are located at the corners of a truncated octahedron with the chloride and four sodium atoms inside. (A similar structure called "carbon sodalite" may occur as a very high pressure form of carbon — see illustration in reference.) Each oxygen atom links between an tetrahedron and an tetrahedron. All the oxygen atoms are equivalent, but one half are in environments that are enantiomorphic to the environments of the other half. The silicon atoms are at the location and symmetry-equivalent positions, and the aluminum ions at the location and symmetry-equivalent positions. The three silicon atoms and the three aluminum atoms listed above closest to a given corner of the unit cell form a six-membered ring of tetrahedra, and the four in any face of the unit cell form a four-membered ring of tetrahedra. The six-membered rings can serve as channels in which ions can diffuse through the crystal.
The structure is a crumpled form of a structure in which the three-fold axes of each tetrahedron lie in planes parallel to the faces of the unit cell, thus putting half the oxygen atoms in the faces. As the temperature is raised the sodalite structure expands and uncrumples, becoming more like this structure. In this structure the two cavities are still chiral, because no indirect isometry centred on the cavity (i.e. a reflexion, inversion, or improper rotation) can superimpose the silicon atoms onto silicon atoms and the aluminum atoms onto aluminum atoms, while also superimposing the sodium atoms on other sodium atoms. A discontinuity of the thermal expansion coefficient occurs at a certain temperature when chloride is replaced by sulfate or iodide, and this is thought to happen when the framework becomes fully expanded or when the cation (sodium in natural sodalite) reaches the coordinates (et cetera). This adds symmetry (such as mirror planes in the faces of the unit cell) so that the space group becomes Pmn (space group 223), and the cavities cease to be chiral and take on pyritohedral symmetry.
Natural sodalite holds primarily chloride anions in the cages, but they can be substituted by other anions such as sulfate, sulfide, hydroxide, trisulfur with other minerals in the sodalite group representing end member compositions. The sodium can be replaced by other alkali group elements, and the chloride by other halides. Many of these have been synthesized.
The characteristic blue color arises mainly from caged and clusters.
Properties
A light, relatively hard yet fragile mineral, sodalite is named after its sodium content; in mineralogy it may be classed as a feldspathoid. Well known for its blue color, sodalite may also be grey, yellow, green, or pink and is often mottled with white veins or patches. The more uniformly blue material is used in jewellery, where it is fashioned into cabochons and beads. Lesser material is more often seen as facing or inlay in various applications.
Although somewhat similar to lazurite and lapis lazuli, sodalite rarely contains pyrite (a common inclusion in lapis) and its blue color is more like traditional royal blue rather than ultramarine. It is further distinguished from similar minerals by its white (rather than blue) streak. Sodalite's six directions of poor cleavage may be seen as incipient cracks running through the stone.
Most sodalite will fluoresce orange under ultraviolet light, and hackmanite exhibits tenebrescence.
Hackmanite
Hackmanite is a variety of sodalite exhibiting tenebrescence. When hackmanite from Mont Saint-Hilaire (Quebec) or Ilímaussaq (Greenland) is freshly quarried, it is generally pale to deep violet but the color fades quickly to greyish or greenish white. Conversely, hackmanite from Afghanistan and the Myanmar Republic starts off creamy white but develops a violet to pink-red color in sunlight. If left in a dark environment for some time, the violet will fade again. Tenebrescence is accelerated by the use of longwave or, particularly, shortwave ultraviolet light.
Occurrence
Sodalite was first described in 1811 for the occurrence in its type locality in the Ilimaussaq complex, Narsaq, West Greenland.
Occurring typically in massive form, sodalite is found as vein fillings in plutonic igneous rocks such as nepheline syenites. It is associated with other minerals typical of silica-undersaturated environments, namely leucite, cancrinite and natrolite. Other associated minerals include nepheline, titanian andradite, aegirine, microcline, sanidine, albite, calcite, fluorite, ankerite and baryte.
Significant deposits of fine material are restricted to but a few locales: Bancroft, Ontario (Princess Sodalite Mine), and Mont-Saint-Hilaire, Quebec, in Canada; and Litchfield, Maine, and Magnet Cove, Arkansas, in the US. The Ice River complex, near Golden, British Columbia, contains sodalite. Smaller deposits are found in South America (Brazil and Bolivia), Portugal, Romania, Burma and Russia. Hackmanite is found principally in Mont-Saint-Hilaire and Greenland.
Euhedral, transparent crystals are found in northern Namibia and in the lavas of Vesuvius, Italy.
Sodalitite is a type of extrusive igneous rock rich in sodalite. Its intrusive equivalent is sodalitolite.
History
The people of the Caral culture traded for sodalite from the Collao altiplano.
Synthesis
The mesoporous cage structure of sodalite makes it useful as a container material for many anions. Some of the anions known to have been included in sodalite-structure materials include nitrate, iodide, iodate, permanganate, perchlorate, and perrhenate.
| Physical sciences | Silicate minerals | Earth science |
456204 | https://en.wikipedia.org/wiki/Organic%20geochemistry | Organic geochemistry | Organic geochemistry is the study of the impacts and processes that organisms have had on the Earth. It is mainly concerned with the composition and mode of origin of organic matter in rocks and in bodies of water. The study of organic geochemistry is traced to the work of Alfred E. Treibs, "the father of organic geochemistry." Treibs first isolated metalloporphyrins from petroleum. This discovery established the biological origin of petroleum, which was previously poorly understood. Metalloporphyrins in general are highly stable organic compounds, and the detailed structures of the extracted derivatives made clear that they originated from chlorophyll.
Applications
Energy
Petroleum
The relationship between the occurrence of organic compounds in sedimentary deposits and petroleum deposits has long been of interest. Studies of ancient sediments and rock provide insights into the origins and sources of oil and petroleum, as well as the biochemical antecedents of life. Oil spills in particular have been of interest to geochemists in regards to the impact of petroleum and oil on the current geological environment. Following the Exxon Valdez Oil Spill, organic geochemistry knowledge on oil-spill chemistry bloomed with the analyses of samples from the spill.
Geochemists study petroleum-inclusions in geological samples to compare present-day fluid-inclusions to dated samples. This analysis provides insight into the age of the petroleum samples and the surrounding rock. Spectrographic, optical, destructive, and nondestructive methods are used to analyze samples via mass spectrometry or Raman spectroscopy. The discovered differences in samples, such as oil-to-gas ratio or viscosity are typically attributed to the rock source of the sample. Other characteristics typically noted are pressure/volume/temperature properties, sample texture, and sample composition. Complications in analysis arise when the source rock is near or in a water source.
Petroleum is also studied via carbon isotope analysis. Carbon isotopes provide insight into the Earth's carbon cycle and geological processes. Geochemists are able to discern the composition of petroleum deposits by examining the ratio of carbon isotopes and comparing this ratio to known values for carbon based structures of which the petroleum could be composed.
Coal
Vast knowledge about coal has been attained since the inception of its use as an energy source. However, modern geochemists are still studying how plant material changes into coal. They have determined coalification results from a selective degradation of plant materials, while other plant material is preserved. Coal macromolecules are usually made up of these degradation-resistant biopolymers contained in algae, spores, and wood. Geochemists have unraveled the mysteries behind coal formation by comparing properties of the biopolymers to properties found in existing coal macromolecules. The analytical methods of Carbon NMR and gas chromatography-mass spectrometry (GC-MS) combined with flash pyrolysis has greatly enhanced the ability of organic geochemists to analyse the minute structural units of coal.
Further knowledge into the age of coal sediments has been attained via isochron dating of uranium in the coalified samples. Examination of the parent to daughter ratio of uranium isotopes has led to the dating of select samples to the Late Cretaceous Period.
Environmental
Modern organic geochemistry includes studies of recent sediments to understand the carbon cycle, climate change, and ocean processes. In connection with petroleum studies, petroleum-focused geochemists also examine the impact of petroleum on the geological environment. Geochemistry also examines other pollutants in geological systems, such as metabolites formed from the degradation of hydrocarbons. Organic geochemistry analytical techniques, such as GC-MS, allow chemists to determine the intricate effects of organic metabolites and human-derived waste products on the geological environment.
Of specific concern are the human-derived pollutants stemming from agricultural work. The use of animal manure, in combination with general municipal and sewage waste management, has changed many physical properties of the agricultural soil involved and the surrounding soils.
Organic geochemistry is also relevant to aqueous environments. Pollutants, their metabolites, and how both enter bodies of water are of particular importance in the field. This organic matter can also be derived from geological processes in or near bodies of water, similarly influencing nearby lifeforms and protein production. Fluorescence spectroscopy has been introduced as a technique to examine organic matter in bodies of water, as dissolved organic matter is typically fluorescent.
The study of organic geochemistry also extends to the atmosphere. Particularly, geochemists in this field study the makeup of insoluble material in the lower atmosphere. They have defined certain consequences of organic aerosols including physiological toxicity, direct and indirect climate forcing, smog, rain acidification, and incorporation into the natural carbon cycle.
| Physical sciences | Geochemistry | Earth science |
456234 | https://en.wikipedia.org/wiki/Colligative%20properties | Colligative properties | In chemistry, colligative properties are those properties of solutions that depend on the ratio of the number of solute particles to the number of solvent particles in a solution, and not on the nature of the chemical species present.<ref>McQuarrie, Donald, et al. Colligative properties of Solutions" General Chemistry Mill Valley: Library of Congress, 2011. .</ref> The number ratio can be related to the various units for concentration of a solution such as molarity, molality, normality (chemistry), etc. The assumption that solution properties are independent of nature of solute particles is exact only for ideal solutions, which are solutions that exhibit thermodynamic properties analogous to those of an ideal gas, and is approximate for dilute real solutions. In other words, colligative properties are a set of solution properties that can be reasonably approximated by the assumption that the solution is ideal.
Only properties which result from the dissolution of a nonvolatile solute in a volatile liquid solvent are considered. They are essentially solvent properties which are changed by the presence of the solute. The solute particles displace some solvent molecules in the liquid phase and thereby reduce the concentration of solvent and increase its entropy, so that the colligative properties are independent of the nature of the solute. The word colligative is derived from the Latin colligatus meaning bound together. This indicates that all colligative properties have a common feature, namely that they are related only to the number of solute molecules relative to the number of solvent molecules and not to the nature of the solute.
Colligative properties include:
Relative lowering of vapor pressure (Raoult's law)
Elevation of boiling point
Depression of freezing point
Osmotic pressure
For a given solute-solvent mass ratio, all colligative properties are inversely proportional to solute molar mass.
Measurement of colligative properties for a dilute solution of a non-ionized solute such as urea or glucose in water or another solvent can lead to determinations of relative molar masses, both for small molecules and for polymers which cannot be studied by other means. Alternatively, measurements for ionized solutes can lead to an estimation of the percentage of dissociation taking place.
Colligative properties are studied mostly for dilute solutions, whose behavior may be approximated as that of an ideal solution. In fact, all of the properties listed above are colligative only in the dilute limit: at higher concentrations, the freezing point depression, boiling point elevation, vapor pressure elevation or depression, and osmotic pressure are all dependent on the chemical nature of the solvent and the solute.
Relative lowering of vapor pressure
A vapor is a substance in a gaseous state at a temperature lower than its critical point. Vapor Pressure is the pressure exerted by a vapor in thermodynamic equilibrium with its solid or liquid state. The vapor pressure of a solvent is lowered when a non-volatile solute is dissolved in it to form a solution.
For an ideal solution, the equilibrium vapor pressure is given by Raoult's law as
where is the vapor pressure of the pure component (i= A, B, ...) and is the mole fraction of the component in the solution.
For a solution with a solvent (A) and one non-volatile solute (B), and .
The vapor pressure lowering relative to pure solvent is , which is proportional to the mole fraction of solute.
If the solute dissociates in solution, then the number of moles of solute is increased by the van 't Hoff factor , which represents the true number of solute particles for each formula unit. For example, the strong electrolyte MgCl2 dissociates into one Mg2+ ion and two Cl− ions, so that if ionization is complete, i = 3 and , where is calculated with moles of solute i times initial moles and moles of solvent same as initial moles of solvent before dissociation. The measured colligative properties show that i is somewhat less than 3 due to ion association.
Boiling point and freezing point
Addition of solute to form a solution stabilizes the solvent in the liquid phase, and lowers the solvent's chemical potential so that solvent molecules have less tendency to move to the gas or solid phases. As a result, liquid solutions slightly above the solvent boiling point at a given pressure become stable, which means that the boiling point increases. Similarly, liquid solutions slightly below the solvent freezing point become stable meaning that the freezing point decreases. Both the boiling point elevation and the freezing point depression are proportional to the lowering of vapor pressure in a dilute solution.
These properties are colligative in systems where the solute is essentially confined to the liquid phase. Boiling point elevation (like vapor pressure lowering) is colligative for non-volatile solutes where the solute presence in the gas phase is negligible. Freezing point depression is colligative for most solutes since very few solutes dissolve appreciably in solid solvents.
Boiling point elevation (ebullioscopy)
The boiling point of a liquid at a given external pressure is the temperature () at which the vapor pressure of the liquid equals the external pressure. The normal boiling point is the boiling point at a pressure equal to 1 atm.
The boiling point of a pure solvent is increased by the addition of a non-volatile solute, and the elevation can be measured by ebullioscopy. It is found that
Here i is the van 't Hoff factor as above, Kb is the ebullioscopic constant of the solvent (equal to 0.512 °C kg/mol for water), and m is the molality of the solution.
The boiling point is the temperature at which there is equilibrium between liquid and gas phases. At the boiling point, the number of gas molecules condensing to liquid equals the number of liquid molecules evaporating to gas. Adding a solute dilutes the concentration of the liquid molecules and reduces the rate of evaporation. To compensate for this and re-attain equilibrium, the boiling point occurs at a higher temperature.
If the solution is assumed to be an ideal solution, Kb can be evaluated from the thermodynamic condition for liquid-vapor equilibrium. At the boiling point, the chemical potential μA of the solvent in the solution phase equals the chemical potential in the pure vapor phase above the solution.
The asterisks indicate pure phases. This leads to the result , where R is the molar gas constant, M is the solvent molar mass and ΔHvap is the solvent molar enthalpy of vaporization.
Freezing point depression (cryoscopy)
The freezing point () of a pure solvent is lowered by the addition of a solute which is insoluble in the solid solvent, and the measurement of this difference is called cryoscopy. It is found that
(which can also be written as )
Here Kf is the cryoscopic constant (equal to 1.86 °C kg/mol for the freezing point of water), i is the van 't Hoff factor, and m the molality (in mol/kg). This predicts the melting of ice by road salt.
In the liquid solution, the solvent is diluted by the addition of a solute, so that fewer molecules are available to freeze. Re-establishment of equilibrium is achieved at a lower temperature at which the rate of freezing becomes equal to the rate of liquefying. At the lower freezing point, the vapor pressure of the liquid is equal to the vapor pressure of the corresponding solid, and the chemical potentials of the two phases are equal as well. The equality of chemical potentials permits the evaluation of the cryoscopic constant as , where ΔfusH is the solvent molar enthalpy of fusion.
Osmotic pressure
The osmotic pressure of a solution is the difference in pressure between the solution and the pure liquid solvent when the two are in equilibrium across a semipermeable membrane, which allows the passage of solvent molecules but not of solute particles. If the two phases are at the same initial pressure, there is a net transfer of solvent across the membrane into the solution known as osmosis. The process stops and equilibrium is attained when the pressure difference equals the osmotic pressure.
Two laws governing the osmotic pressure of a dilute solution were discovered by the German botanist W. F. P. Pfeffer and the Dutch chemist J. H. van’t Hoff:
The osmotic pressure of a dilute solution at constant temperature is directly proportional to its concentration.
The osmotic pressure of a solution is directly proportional to its absolute temperature.
These are analogous to Boyle's law and Charles's law for gases. Similarly, the combined ideal gas law, , has as an analogue for ideal solutions , where is osmotic pressure; V is the volume; n is the number of moles of solute; R is the molar gas constant 8.314 J K−1 mol−1; T is absolute temperature; and i is the Van 't Hoff factor.
The osmotic pressure is then proportional to the molar concentration , since
The osmotic pressure is proportional to the concentration of solute particles ci and is therefore a colligative property.
As with the other colligative properties, this equation is a consequence of the equality of solvent chemical potentials of the two phases in equilibrium. In this case the phases are the pure solvent at pressure P and the solution at total pressure (P + ).
History
The word colligative (Latin: co, ligare) was introduced in 1891 by Wilhelm Ostwald. Ostwald classified solute properties in three categories:H.W. Smith, Circulation 21, 808 (1960) Theory of Solutions: A Knowledge of the Laws of Solutions ...
colligative properties, which depend only on solute concentration and temperature and are independent of the nature of the solute particles
additive properties such as mass, which are the sums of properties of the constituent particles and therefore depend also on the composition (or molecular formula) of the solute, and
constitutional'' properties, which depend further on the molecular structure of the given solute.
| Physical sciences | Thermodynamics | Chemistry |
456371 | https://en.wikipedia.org/wiki/Fire%20salamander | Fire salamander | The fire salamander (Salamandra salamandra) is a common species of salamander found in Europe.
It is black with yellow spots or stripes to a varying degree; some specimens can be nearly completely black while on others the yellow is dominant. Shades of red and orange may sometimes appear, either replacing or mixing with the yellow according to subspecies. This bright coloration is highly conspicuous and acts to deter predators by honest signalling of its toxicity (aposematism). Fire salamanders can have a very long lifespan; one specimen lived for more than 50 years in Museum Koenig, a German natural history museum.
Despite its wide distribution and abundance, it is classified as Vulnerable on the IUCN Red List due to its susceptibility to infection by the introduced fungus Batrachochytrium salamandrivorans, which has caused severe declines in fire salamanders in parts of its range.
Taxonomy
Several subspecies of the fire salamander are recognized. Most notable are the subspecies fastuosa and bernadezi, which are the only viviparous subspecies – the others are ovoviviparous.
S. s. alfredschmidti
S. s. almanzoris
S. s. bejarae
S. s. bernardezi
S. s. beschkovi
S. s. crespoi
S. s. fastuosa (or bonalli) – yellow-striped fire salamander
S. s. gallaica – Galician fire salamander
S. s. gigliolii
S. s. morenica
S. s. salamandra – spotted fire salamander, nominate subspecies
S. s. terrestris – barred fire salamander
S. s. werneri
Some former subspecies have been lately recognized as species for genetic reasons.
S. algira Bedriaga, 1883 – African fire salamander
S. corsica Savi, 1838 – Corsican fire salamander
S. infraimmaculata Martens, 1885 – Near Eastern fire salamander (arouss al ayn)
Distribution
Fire salamanders are found in most of southern and central Europe. They are most commonly found at altitudes between and , only rarely below (in Northern Germany sporadically down to ). However, in the Balkans or Spain they are commonly found in higher altitudes as well.
The scientific article titled "Water, Stream Morphology and Landscape: Complex Habitat Determinants for the Fire Salamander Salamandra salamandra" explored the factors influencing the distribution of the fire salamander, a semiaquatic amphibian species, in northern Italy. The study aimed to understand the relationship between environmental features and species distribution, essential for effective habitat conservation.
Researchers evaluated three main factors: stream morphology, biotic features of water, and the composition of the surrounding landscape near wetlands. They collected data from 132 localities over four years and used an information-theoretic approach to build species distribution models. Variance partitioning was then employed to assess the relative importance of environmental variables.
The findings revealed that the distribution of fire salamander larvae was associated with specific environmental conditions. They were found in heterogeneous and shallow streams with scarce periphyton (a type of algae) and rich macrobenthos (aquatic invertebrates), characteristic of oligotrophic water. Additionally, the presence of woodlands in the surrounding landscape played a crucial role in the species' distribution.
The study emphasized the interconnectedness of multiple factors in determining Salamandra salamandra distribution. Stream morphology was the most influential variable, but the combined effects of water features and landscape composition also played significant roles. The article underscores the importance of considering both aquatic and upland habitats in conservation efforts for these and other semiaquatic amphibians.
Genetic differentiation by population
A 2021 research project investigated the role of physical and ecological isolation in shaping genetic differentiation patterns among populations and subspecies of the fire salamander in central Iberia. Researchers utilized microsatellite genetic data and environmental dissimilarity measures to assess the impact of both types of isolation on genetic connectivity.
The analysis revealed significant genetic diversity variation across the study area, with lower diversity in eastern populations near the range limit and higher diversity in western and central populations. The study identified strong genetic structure, as populations from the Iberian Central System (ICS) and the Montes de Toledo Range (MTR) formed distinct genetic groups. Physical isolation, represented by landscape resistance, played a substantial role in genetic differentiation between populations across all spatial extents. Different types of landscape resistance, such as climate-based and landcover-based, provided the best model fits in different regions. The researchers proposed a scenario where gene flow between two subspecies, S. s. bejarae and S. s. almanzoris, was restricted by ecological isolation associated with sharp transitions in precipitation seasonality. However, gene flow between populations with intermediate levels of precipitation seasonality was less restricted. The results provided evidence for ongoing environmental adaptation, leading to the maintenance of distinct ecotypes and evolutionary units.
Habitat, behavior and diet
Fire salamanders live in the forests of central Europe and are more common in hilly areas. They prefer deciduous forests since they like to hide in fallen leaves and around mossy tree trunks. They need small brooks or ponds with clean water in their habitat for the development of the larvae. Whether on land or in water, fire salamanders are inconspicuous. They spend much of their time hidden under wood or other objects. They are active in the evening and the night, but on rainy days they are active in the daytime as well.
The diet of the fire salamander consists of various insects, spiders, millipedes, centipedes, earthworms and slugs, but they also occasionally eat newts and young frogs. In captivity, they eat crickets, mealworms, waxworms and silkworm larvae. Small prey will be caught within the range of the vomerine teeth or by the posterior half of the tongue, to which the prey adheres. It weighs about 40 grams. Compared to other salamanders in the region like Luschan's salamander, the fire salamander has been shown to be larger and appears to have a more solid pectoral girdle. Additionally, it has a longer pectoral girdle than Luschan’s salamander. The fire salamander is one of Europe's largest salamanders and can grow to be long.
Diet and habitat interaction
A study in 2013 aimed to investigate the foraging behavior of fire salamander larvae from different environments, specifically caves and streams, and to understand the roles of local adaptation and phenotypic plasticity in shaping their behavior. The researchers conducted a behavioral experiment using newborn larvae from 11 caves and nine streams in northwest Italy. In the experiment, the larvae were individually maintained in laboratory conditions and subjected to different test conditions, including light/darkness, prey presence/absence, and food deprivation/normal feeding. Video tracking was used to quantify the larvae's movements and foraging strategies.
The results revealed significant differences in foraging behavior between cave and stream larvae. The cave larvae exhibited a more active foraging strategy, especially in darkness and in the absence of prey, suggesting local adaptations to the challenging cave environment with limited food resources. Stream larvae, on the other hand, preferred using peripheral sectors of the test arena, indicating a preference for sit-and-wait behavior, which is advantageous in the presence of detectable and active prey.
The study demonstrated that fire salamander larvae are highly plastic in their foraging behavior. They adjusted their activity levels and movement patterns in response to changes in light conditions, prey availability, and food deprivation. The plastic responses observed were beneficial for increasing encounter rates with prey and optimizing energy utilization in resource-scarce environments. The study revealed an interplay between phenotypic plasticity and local adaptation in shaping the foraging behavior of fire salamander larvae. While plasticity appears to be dominant in the early stages of colonization and adaptation to new environments, local adaptations may also contribute to behavioral differences between cave and stream populations.
Reproduction
Males and females look very similar, except during the breeding season, when the most conspicuous difference is a swollen gland around the male's vent. This gland produces the spermatophore, which carries a sperm packet at its tip. The courtship happens on land. After the male becomes aware of a potential mate, he confronts her and blocks her path. The male rubs her with his chin to express his interest in mating, then crawls beneath her and grasps her front limbs with his own in amplexus. He deposits a spermatophore on the ground, then attempts to lower the female's cloaca into contact with it. If successful, the female draws the sperm packet in and her eggs are fertilized internally. The eggs develop internally and the female deposits the larvae into a body of water just as they hatch. In some subspecies, the larvae continue to develop within the female until she gives birth to fully formed metamorphs. Breeding has not been observed in neotenic fire salamanders.
In captivity, females may retain sperm long-term and use the stored sperm later to produce another clutch. This behavior has not been observed in the wild, likely due to the ability to obtain fresh sperm and the degradation of stored sperm.
Experimental and cave reproduction
A European study investigated the breeding and developmental patterns of the fire salamander in both natural and artificial caves across various regions in Italy. The researchers conducted extensive surveys from 2008 to 2017, exploring a total of 292 sites, comprising 219 natural caves and 73 artificial caves. Among these sites, 52 were found to host underground breeding sites of fire salamanders, with 15 occurring in natural caves and 37 in artificial sites.
The experiment explored environmental features in determining larval distribution inside caves. Fire salamander larvae were observed to choose caves with specific characteristics, such as stable water presence, ease of access, and the presence of rich macrobenthos communities. Larval development in underground springs and natural caves was found to be slower compared to epigean environments, possibly influenced by factors such as temperature and food availability. Furthermore, the lack of light in caves influenced the predation behavior of larvae, with cave populations showing higher adaptability in capturing prey.
Cave environments presented unique challenges for fire salamanders, including food scarcity and the occurrence of cannibalism, particularly in resource-poor habitats. However, the study revealed that fire salamanders exhibited strong phenotypic plasticity, which allowed them to adapt and survive in these extreme underground conditions.
The research emphasizes the importance of local adaptations and phenotypic plasticity in the successful colonization of caves by fire salamanders. It also highlights the need for further genetic studies to understand the differentiation between cave and stream populations and the mechanisms driving successful cave exploitation. Despite challenges posed by large urodele genomes, future genome scan and transcriptomic approaches may provide valuable insights into the genetic processes involved in cave adaptation.
Toxicity
The fire salamander's primary alkaloid toxin, samandarin, causes strong muscle convulsions and hypertension combined with hyperventilation in all vertebrates. Through an analysis of the European fire salamander’s skin secretions, scientists have determined that another alkaloid, such as samandarone, is also released by the salamander. These steroids can be swabbed from the salamander’s parotid glands. Samandarine was often the dominant alkaloid present but the ratio varied between salamanders. This ratio, however, was not shown to be sex dependent. Larvae do not produce these alkaloids. Upon maturity, ovaries, livers, and testes appear to produce these defensive steroids.
The poison glands of the fire salamander are concentrated in certain areas of the body, especially around the head and the dorsal skin surface. The coloured portions of the animal's skin usually coincide with these glands. Compounds in the skin secretions may be effective against bacterial and fungal infections of the epidermis; some are potentially dangerous to human life.
A 2002 study focused on investigating the variability of toxic alkaloids in the skin secretion of the European fire salamander. The chemical defense mechanisms of the salamander provides valuable insights into the chemical composition of skin secretions in amphibians. The two major alkaloids of focus were, samandarine and samandarone. Using gas chromatography/mass spectrometry, the researchers analyzed individual specimens from two populations of fire salamanders and observed a high degree of intraspecific variability in the ratio of samandarine to samandarone in the skin secretion. Some individuals had a higher concentration of samandarone, while others exhibited equal levels of both alkaloids.
Internal organs contained either no or only small amounts of the alkaloids, and the ratio of alkaloids in the organs differed from that in the skin. Particularly noteworthy was the finding that the larvae found in the oviducts of gravid females were entirely free of alkaloids, and their skin lacked the typical granular glands that are present in adult salamanders. Samandarone may be a product of a separate biosynthetic pathway due to its exclusive presence in skin secretions and organ extracts.
Environmental stressors and threats
Introduced species
A research team conducted a comprehensive study investigating the impact of mosquitofish (Gambusia affinis) on endangered Near Eastern fire salamander (S. infraimmaculata) larvae in Israel. The research was conducted through a combination of field surveys and a mesocosm experiment to understand the potential threat posed by mosquitofish to the native amphibian population.
Researchers observed natural breeding pools of fire salamanders, both with and without mosquitofish. The presence of mosquitofish was found to have a detrimental effect on the salamander larvae, leading to reduced densities, smaller sizes, and lower tail:body ratios in the pools with mosquitofish. These observations indicated that mosquitofish predation was causing severe physical damage to the salamander larvae.
To further investigate, a mesocosm experiment was conducted. The researchers manipulated the presence of mosquitofish and structural complexity in the artificial breeding pools. The results supported the field observations, showing that mosquitofish had a significant negative impact on salamander survival, size, and body condition. The fish-inflicted damage included partial tail fins, gill injuries, and limb damage, leading to a reduced likelihood of successful metamorphosis for the salamander larvae in mosquitofish-present mesocosms.
Importantly, the study revealed that increased structural complexity (artificial vegetation) did not provide a refuge for the salamander larvae against mosquitofish predation, contrary to expectations. It was also noted that the use of mosquitofish for mosquito control in permanent ponds could lead to negative consequences for native amphibian populations, as the presence of mosquitofish posed a significant threat to the survival of the fire salamander larvae.
The experiment suggests that mosquitofish pose a serious threat to the endangered fire salamander population in Israel. It also highlights conserving the native amphibian species by reconsidering the use of mosquitofish for mosquito control in habitats where these vulnerable species breed. Efforts to remove mosquitofish from Salamandra-breeding sites are recommended to safeguard the long-term persistence of the fire salamander population and protect against potential ecological disruptions caused by invasive fish species.
Batrachochytrium salamandrivorans
In parts of its range, the fire salamander has become highly endangered by the spread of the introduced chytrid fungus Batrachochytrium salamandrivorans, which has had catastrophic effects on its population. This collapse was first identified from the Netherlands in 2013. The fire salamander in the Netherlands is teetering on the brink of extinction, confined to three small populations in the southern part of the country. Prior to these declines, they were already listed as "Endangered" on the national Red List, and their range had reduced by 57% since 1950, mainly due to changes in water availability and habitat degradation. The remaining populations were limited to specific areas of deciduous forests on hillsides, and their surface activity is restricted to humid periods with night temperatures above 5°C.The species had been considered stable until 2008 when dead individuals were observed, and since 2010, there has been a staggering 96% population decline, with the largest population dropping from 241 individuals to only four in 2011. In 2013, the cause of the decline was officially identified as a new chytrid fungus, Batrachochytrium salamandrivorans (Bsal), likely introduced to Europe from east Asia via captive amphibians.
Since its identification in the Netherlands, Bsal has continued to spread across western Europe, and has infected more populations of S. s. terrestris in Belgium and western Germany, with an isolated but contained occurrence in Spain affecting a population of S. s. hispanica. Dramatic declines have been noted in all affected populations, and some may eventually be entirely extirpated, although at most known sites, fire salamanders persist at low numbers even after disease outbreak, and in one case appear to have recovered. Some localities in the Eifel Mountains where fire salamanders were previously known from appear to now be devoid of fire salamanders, suggesting landscape-scale declines that occurred prior to the disease's identification by science. In 2023, the fire salamander was officially moved from 'Least Concern' to 'Vulnerable' on the IUCN Red List, relating to the past and predicted future declines in the species.
Gallery
| Biology and health sciences | Salamanders and newts | Animals |
456410 | https://en.wikipedia.org/wiki/Crystal%20system | Crystal system | In crystallography, a crystal system is a set of point groups (a group of geometric symmetries with at least one fixed point). A lattice system is a set of Bravais lattices. Space groups are classified into crystal systems according to their point groups, and into lattice systems according to their Bravais lattices. Crystal systems that have space groups assigned to a common lattice system are combined into a crystal family.
The seven crystal systems are triclinic, monoclinic, orthorhombic, tetragonal, trigonal, hexagonal, and cubic. Informally, two crystals are in the same crystal system if they have similar symmetries (though there are many exceptions).
Classifications
Crystals can be classified in three ways: lattice systems, crystal systems and crystal families. The various classifications are often confused: in particular the trigonal crystal system is often confused with the rhombohedral lattice system, and the term "crystal system" is sometimes used to mean "lattice system" or "crystal family".
Lattice system
A lattice system is a group of lattices with the same set of lattice point groups. The 14 Bravais lattices are grouped into seven lattice systems: triclinic, monoclinic, orthorhombic, tetragonal, rhombohedral, hexagonal, and cubic.
Crystal system
A crystal system is a set of point groups in which the point groups themselves and their corresponding space groups are assigned to a lattice system. Of the 32 crystallographic point groups that exist in three dimensions, most are assigned to only one lattice system, in which case both the crystal and lattice systems have the same name. However, five point groups are assigned to two lattice systems, rhombohedral and hexagonal, because both exhibit threefold rotational symmetry. These point groups are assigned to the trigonal crystal system.
Crystal family
A crystal family is determined by lattices and point groups. It is formed by combining crystal systems that have space groups assigned to a common lattice system. In three dimensions, the hexagonal and trigonal crystal systems are combined into one hexagonal crystal family.
Comparison
Five of the crystal systems are essentially the same as five of the lattice systems. The hexagonal and trigonal crystal systems differ from the hexagonal and rhombohedral lattice systems. These are combined into the hexagonal crystal family.
The relation between three-dimensional crystal families, crystal systems and lattice systems is shown in the following table:
Note: there is no "trigonal" lattice system. To avoid confusion of terminology, the term "trigonal lattice" is not used.
Crystal classes
The 7 crystal systems consist of 32 crystal classes (corresponding to the 32 crystallographic point groups) as shown in the following table below:
The point symmetry of a structure can be further described as follows. Consider the points that make up the structure, and reflect them all through a single point, so that (x,y,z) becomes (−x,−y,−z). This is the 'inverted structure'. If the original structure and inverted structure are identical, then the structure is centrosymmetric. Otherwise it is non-centrosymmetric. Still, even in the non-centrosymmetric case, the inverted structure can in some cases be rotated to align with the original structure. This is a non-centrosymmetric achiral structure. If the inverted structure cannot be rotated to align with the original structure, then the structure is chiral or enantiomorphic and its symmetry group is enantiomorphic.
A direction (meaning a line without an arrow) is called polar if its two-directional senses are geometrically or physically different. A symmetry direction of a crystal that is polar is called a polar axis. Groups containing a polar axis are called polar. A polar crystal possesses a unique polar axis (more precisely, all polar axes are parallel). Some geometrical or physical property is different at the two ends of this axis: for example, there might develop a dielectric polarization as in pyroelectric crystals. A polar axis can occur only in non-centrosymmetric structures. There cannot be a mirror plane or twofold axis perpendicular to the polar axis, because they would make the two directions of the axis equivalent.
The crystal structures of chiral biological molecules (such as protein structures) can only occur in the 65 enantiomorphic space groups (biological molecules are usually chiral).
Bravais lattices
There are seven different kinds of lattice systems, and each kind of lattice system has four different kinds of centerings (primitive, base-centered, body-centered, face-centered). However, not all of the combinations are unique; some of the combinations are equivalent while other combinations are not possible due to symmetry reasons. This reduces the number of unique lattices to the 14 Bravais lattices.
The distribution of the 14 Bravais lattices into 7 lattice systems is given in the following table.
In geometry and crystallography, a Bravais lattice is a category of translative symmetry groups (also known as lattices) in three directions.
Such symmetry groups consist of translations by vectors of the form
R = n1a1 + n2a2 + n3a3,
where n1, n2, and n3 are integers and a1, a2, and a3 are three non-coplanar vectors, called primitive vectors.
These lattices are classified by the space group of the lattice itself, viewed as a collection of points; there are 14 Bravais lattices in three dimensions; each belongs to one lattice system only. They represent the maximum symmetry a structure with the given translational symmetry can have.
All crystalline materials (not including quasicrystals) must, by definition, fit into one of these arrangements.
For convenience a Bravais lattice is depicted by a unit cell which is a factor 1, 2, 3, or 4 larger than the primitive cell. Depending on the symmetry of a crystal or other pattern, the fundamental domain is again smaller, up to a factor 48.
The Bravais lattices were studied by Moritz Ludwig Frankenheim in 1842, who found that there were 15 Bravais lattices. This was corrected to 14 by A. Bravais in 1848.
In other dimensions
Two-dimensional space
In two-dimensional space, there are four crystal systems (oblique, rectangular, square, hexagonal), four crystal families (oblique, rectanguar, square, hexagonal), and four lattice systems (oblique, rectangular, square, and hexagonal).
Four-dimensional space
The four-dimensional unit cell is defined by four edge lengths (a, b, c, d) and six interaxial angles (α, β, γ, δ, ε, ζ). The following conditions for the lattice parameters define 23 crystal families
The names here are given according to Whittaker. They are almost the same as in Brown et al., with exception for names of the crystal families 9, 13, and 22. The names for these three families according to Brown et al. are given in parentheses.
The relation between four-dimensional crystal families, crystal systems, and lattice systems is shown in the following table. Enantiomorphic systems are marked with an asterisk. The number of enantiomorphic pairs is given in parentheses. Here the term "enantiomorphic" has a different meaning than in the table for three-dimensional crystal classes. The latter means, that enantiomorphic point groups describe chiral (enantiomorphic) structures. In the current table, "enantiomorphic" means that a group itself (considered as a geometric object) is enantiomorphic, like enantiomorphic pairs of three-dimensional space groups P31 and P32, P4122 and P4322. Starting from four-dimensional space, point groups also can be enantiomorphic in this sense.
| Physical sciences | Crystallography | Physics |
456729 | https://en.wikipedia.org/wiki/Computer-aided%20engineering | Computer-aided engineering | Computer-aided engineering (CAE) is the general usage of technology to aid in tasks related to engineering analysis. Any use of technology to solve or assist engineering issues falls under this umbrella.
Overview
Following alongside the consistent improvement in computer graphics and speed, computer aid assists engineers with once complicated and time consuming tasks with the input of information and a press of a button.
It includes finite element method or analysis (FEA), computational fluid dynamics (CFD), multibody dynamics (MBD), durability and optimization. It is included with computer-aided design (CAD) and computer-aided manufacturing (CAM) in a collective term and abbreviation computer-aided technologies (CAx).
The term CAE has been used to describe the use of computer technology within engineering in a broader sense than just engineering analysis. It was in this context that the term was coined by Jason Lemon, founder of Structural Dynamics Research Corporation (SDRC) in the late 1970s. However, this definition is better known today by the terms CAx and product lifecycle management (PLM).
CAE systems are individually considered a single node on a total information network and each node may interact with other nodes on the network.
CAE fields and phases
CAE areas covered include:
Stress analysis on components and assemblies using finite element analysis (FEA);
Thermal and fluid flow analysis computational fluid dynamics (CFD);
Multibody dynamics (MBD) and kinematics;
Analysis tools for process simulation for operations such as casting, molding, and die press forming;
Optimization of the product or process.
In general, there are three phases in any computer-aided engineering task:
Pre-processing – defining the model and environmental factors to be applied to it (typically a finite element model, but facet, voxel, and thin sheet methods are also used);
Analysis solver (usually performed on high powered computers);
Post-processing of results (using visualization tools).
This cycle is iterated, often many times, either manually or with the use of commercial optimization software.
CAE in the automotive industry
CAE tools are widely used in the automotive industry. Their use has enabled automakers to reduce product development costs and time while improving the safety, comfort, and durability of the vehicles they produce. The predictive capability of CAE tools has progressed to the point where much of the design verification is done using computer simulations (diagnosis) rather than physical prototype testing. CAE dependability is based upon all proper assumptions as inputs and must identify critical inputs (BJ). Even though there have been many advances in CAE, and it is widely used in the engineering field, physical testing is still a must. It is used for verification and model updating, to accurately define loads and boundary conditions, and for final prototype sign-off.
The future of CAE in the product development process
Even though CAE has built a strong reputation as a verification, troubleshooting and analysis tool, there is still a perception that sufficiently accurate results come rather late in the design cycle to really drive the design. This can be expected to become a problem as modern products become ever more complex. They include smart systems, which leads to an increased need for multi-physics analysis including controls, and contain new lightweight materials, with which engineers are often less familiar.
CAE software companies and manufacturers are constantly looking for tools and process improvements to change this situation.
On the software side, they are constantly looking to develop more powerful solvers, to better utilize computer resources, and to include engineering knowledge in pre and post-processing. On the process side, they try to achieve a better alignment between 3D CAE, 1D system simulation, and physical testing. This should increase modeling realism and calculation speed.
CAE software companies and manufacturers try to better integrate CAE in the overall product lifecycle management. In this way they can connect product design with product use, which is needed for smart products. This enhanced engineering process is also referred to as predictive engineering analytics.
| Technology | Disciplines | null |
457036 | https://en.wikipedia.org/wiki/Sulfur%20hexafluoride | Sulfur hexafluoride | Sulfur hexafluoride or sulphur hexafluoride (British spelling) is an inorganic compound with the formula SF6. It is a colorless, odorless, non-flammable, and non-toxic gas. has an octahedral geometry, consisting of six fluorine atoms attached to a central sulfur atom. It is a hypervalent molecule.
Typical for a nonpolar gas, is poorly soluble in water but quite soluble in nonpolar organic solvents. It has a density of 6.12 g/L at sea level conditions, considerably higher than the density of air (1.225 g/L). It is generally stored and transported as a liquefied compressed gas.
has 23,500 times greater global warming potential (GWP) than as a greenhouse gas (over a 100-year time-frame) but exists in relatively minor concentrations in the atmosphere. Its concentration in Earth's troposphere reached 11.50 parts per trillion (ppt) in October 2023, rising at 0.37 ppt/year. The increase since 1980 is driven in large part by the expanding electric power sector, including fugitive emissions from banks of gas contained in its medium- and high-voltage switchgear. Uses in magnesium, aluminium, and electronics manufacturing also hastened atmospheric growth. The 1997 Kyoto Protocol, which came into force in 2005, is supposed to limit emissions of this gas. In a somewhat nebulous way it has been included as part of the carbon emission trading scheme. In some countries this has led to the defunction of entire industries.
Synthesis and reactions
Sulfur hexafluoride on Earth exists primarily as a synthetic industrial gas, but has also been found to occur naturally.
can be prepared from the elements through exposure of to . This was the method used by the discoverers Henri Moissan and Paul Lebeau in 1901. Some other sulfur fluorides are cogenerated, but these are removed by heating the mixture to disproportionate any (which is highly toxic) and then scrubbing the product with NaOH to destroy remaining .
Alternatively, using bromine, sulfur hexafluoride can be synthesized from SF4 and CoF3 at lower temperatures (e.g. 100 °C), as follows:
There is virtually no reaction chemistry for . A main contribution to the inertness of SF6 is the steric hindrance of the sulfur atom, whereas its heavier group 16 counterparts, such as SeF6 are more reactive than SF6 as a result of less steric hindrance. It does not react with molten sodium below its boiling point, but reacts exothermically with lithium. As a result of its inertness, has an atmospheric lifetime of around 3200 years, and no significant environmental sinks other than the ocean.
Applications
By 2000, the electrical power industry is estimated to use about 80% of the sulfur hexafluoride produced, mostly as a gaseous dielectric medium. Other main uses as of 2015 included a silicon etchant for semiconductor manufacturing, and an inert gas for the casting of magnesium.
Dielectric medium
is used in the electrical industry as a gaseous dielectric medium for high-voltage sulfur hexafluoride circuit breakers, switchgear, and other electrical equipment, often replacing oil-filled circuit breakers (OCBs) that can contain harmful polychlorinated biphenyls (PCBs). gas under pressure is used as an insulator in gas insulated switchgear (GIS) because it has a much higher dielectric strength than air or dry nitrogen. The high dielectric strength is a result of the gas's high electronegativity and density. This property makes it possible to significantly reduce the size of electrical gear. This makes GIS more suitable for certain purposes such as indoor placement, as opposed to air-insulated electrical gear, which takes up considerably more room.
Gas-insulated electrical gear is also more resistant to the effects of pollution and climate, as well as being more reliable in long-term operation because of its controlled operating environment. Exposure to an arc chemically breaks down though most of the decomposition products tend to quickly re-form , a process termed "self-healing". Arcing or corona can produce disulfur decafluoride (), a highly toxic gas, with toxicity similar to phosgene. was considered a potential chemical warfare agent in World War II because it does not produce lacrimation or skin irritation, thus providing little warning of exposure.
is also commonly encountered as a high voltage dielectric in the high voltage supplies of particle accelerators, such as Van de Graaff generators and Pelletrons and high voltage transmission electron microscopes.
Alternatives to as a dielectric gas include several fluoroketones. Compact GIS technology that combines vacuum switching with clean air insulation has been introduced for a subset of applications up to 420 kV.
Medical use
is used to provide a tamponade or plug of a retinal hole in retinal detachment repair operations in the form of a gas bubble. It is inert in the vitreous chamber. The bubble initially doubles its volume in 36 hours due to oxygen and nitrogen entering it, before being absorbed in the blood in 10–14 days.
is used as a contrast agent for ultrasound imaging. Sulfur hexafluoride microbubbles are administered in solution through injection into a peripheral vein. These microbubbles enhance the visibility of blood vessels to ultrasound. This application has been used to examine the vascularity of tumours. It remains visible in the blood for 3 to 8 minutes, and is exhaled by the lungs.
Tracer compound
Sulfur hexafluoride was the tracer gas used in the first roadway air dispersion model calibration; this research program was sponsored by the U.S. Environmental Protection Agency and conducted in Sunnyvale, California on U.S. Highway 101. Gaseous is used as a tracer gas in short-term experiments of ventilation efficiency in buildings and indoor enclosures, and for determining infiltration rates. Two major factors recommend its use: its concentration can be measured with satisfactory accuracy at very low concentrations, and the Earth's atmosphere has a negligible concentration of .
Sulfur hexafluoride was used as a non-toxic test gas in an experiment at St John's Wood tube station in London, United Kingdom on 25 March 2007. The gas was released throughout the station, and monitored as it drifted around. The purpose of the experiment, which had been announced earlier in March by the Secretary of State for Transport Douglas Alexander, was to investigate how toxic gas might spread throughout London Underground stations and buildings during a terrorist attack.
Sulfur hexafluoride is also routinely used as a tracer gas in laboratory fume hood containment testing. The gas is used in the final stage of ASHRAE 110 fume hood qualification. A plume of gas is generated inside of the fume hood and a battery of tests are performed while a gas analyzer arranged outside of the hood samples for SF6 to verify the containment properties of the fume hood.
It has been used successfully as a tracer in oceanography to study diapycnal mixing and air-sea gas exchange.
Other uses
The magnesium industry uses as an inert "cover gas" to prevent oxidation during casting, and other processes including smelting. Once the largest user, consumption has declined greatly with capture and recycling.
Insulated glazing windows have used it as a filler to improve their thermal and acoustic insulation performance.
plasma is used in the semiconductor industry as an etchant in processes such as deep reactive-ion etching. A small fraction of the breaks down in the plasma into sulfur and fluorine, with the fluorine ions performing a chemical reaction with silicon.
Tires filled with it take longer to deflate from diffusion through rubber due to the larger molecule size.
Nike likewise used it to obtain a patent and to fill the cushion bags in all of their "Air"-branded shoes from 1992 to 2006. 277 tons was used during the peak in 1997.
The United States Navy's Mark 50 torpedo closed Rankine-cycle propulsion system is powered by sulfur hexafluoride in an exothermic reaction with solid lithium.
Waveguides in high-power microwave systems are pressurized with it. The gas electrically insulates the waveguide, preventing internal arcing.
Electrostatic loudspeakers have used it because of its high dielectric strength and high molecular weight.
Disulfur decafluoride, a chemical weapon, is produced with it as a feedstock.
For entertainment purposes, when breathed, causes the voice to become significantly deeper, due to its density being so much higher than air. This phenomenon is related to the more well-known effect of breathing low-density helium, which causes someone's voice to become much higher. Both of these effects should only be attempted with caution as these gases displace oxygen that the lungs are attempting to extract from the air. Sulfur hexafluoride is also mildly anesthetic.
For science demonstrations / magic as "invisible water" since a light foil boat can be floated in a tank, as will an air-filled balloon.
It is used for benchmark and calibration measurements in Associative and Dissociative Electron Attachment (DEA) experiments
Greenhouse gas
According to the Intergovernmental Panel on Climate Change, is the most potent greenhouse gas. Its global warming potential of 23,900 times that of when compared over a 100-year period. Sulfur hexafluoride is inert in the troposphere and stratosphere and is extremely long-lived, with an estimated atmospheric lifetime of 800–3,200 years.
Measurements of SF6 show that its global average mixing ratio has increased from a steady base of about 54 parts per quadrillion prior to industrialization, to over 11.5 parts per trillion (ppt) as of October 2023, and is increasing by about 0.4 ppt (3.5%) per year. Average global SF6 concentrations increased by about 7% per year during the 1980s and 1990s, mostly as the result of its use in magnesium production, and by electrical utilities and electronics manufacturers. Given the small amounts of SF6 released compared to carbon dioxide, its overall individual contribution to global warming is estimated to be less than 0.2%, however the collective contribution of it and similar man-made halogenated gases has reached about 10% as of 2020. Alternatives are being tested.
In Europe, falls under the F-Gas directive which ban or control its use for several applications. Since 1 January 2006, is banned as a tracer gas and in all applications except high-voltage switchgear. It was reported in 2013 that a three-year effort by the United States Department of Energy to identify and fix leaks at its laboratories in the United States such as the Princeton Plasma Physics Laboratory, where the gas is used as a high voltage insulator, had been productive, cutting annual leaks by . This was done by comparing purchases with inventory, assuming the difference was leaked, then locating and fixing the leaks.
Physiological effects and precautions
Sulfur hexafluoride is a nontoxic gas, but by displacing oxygen in the lungs, it also carries the risk of asphyxia if too much is inhaled. Since it is more dense than air, a substantial quantity of gas, when released, will settle in low-lying areas and present a significant risk of asphyxiation if the area is entered. That is particularly relevant to its use as an insulator in electrical equipment since workers may be in trenches or pits below equipment containing .
As with all gases, the density of affects the resonance frequencies of the vocal tract, thus changing drastically the vocal sound qualities, or timbre, of those who inhale it. It does not affect the vibrations of the vocal folds. The density of sulfur hexafluoride is relatively high at room temperature and pressure due to the gas's large molar mass. Unlike helium, which has a molar mass of about 4 g/mol and pitches the voice up, has a molar mass of about 146 g/mol, and the speed of sound through the gas is about 134 m/s at room temperature, pitching the voice down. For comparison, the molar mass of air, which is about 80% nitrogen and 20% oxygen, is approximately 30 g/mol which leads to a speed of sound of 343 m/s.
Sulfur hexafluoride has an anesthetic potency slightly lower than nitrous oxide; it is classified as a mild anesthetic.
| Physical sciences | Halide salts | Chemistry |
457210 | https://en.wikipedia.org/wiki/Pure%20mathematics | Pure mathematics | Pure mathematics is the study of mathematical concepts independently of any application outside mathematics. These concepts may originate in real-world concerns, and the results obtained may later turn out to be useful for practical applications, but pure mathematicians are not primarily motivated by such applications. Instead, the appeal is attributed to the intellectual challenge and aesthetic beauty of working out the logical consequences of basic principles.
While pure mathematics has existed as an activity since at least ancient Greece, the concept was elaborated upon around the year 1900, after the introduction of theories with counter-intuitive properties (such as non-Euclidean geometries and Cantor's theory of infinite sets), and the discovery of apparent paradoxes (such as continuous functions that are nowhere differentiable, and Russell's paradox). This introduced the need to renew the concept of mathematical rigor and rewrite all mathematics accordingly, with a systematic use of axiomatic methods. This led many mathematicians to focus on mathematics for its own sake, that is, pure mathematics.
Nevertheless, almost all mathematical theories remained motivated by problems coming from the real world or from less abstract mathematical theories. Also, many mathematical theories, which had seemed to be totally pure mathematics, were eventually used in applied areas, mainly physics and computer science. A famous early example is Isaac Newton's demonstration that his law of universal gravitation implied that planets move in orbits that are conic sections, geometrical curves that had been studied in antiquity by Apollonius. Another example is the problem of factoring large integers, which is the basis of the RSA cryptosystem, widely used to secure internet communications.
It follows that, presently, the distinction between pure and applied mathematics is more a philosophical point of view or a mathematician's preference rather than a rigid subdivision of mathematics.
History
Ancient Greece
Ancient Greek mathematicians were among the earliest to make a distinction between pure and applied mathematics. Plato helped to create the gap between "arithmetic", now called number theory, and "logistic", now called arithmetic. Plato regarded logistic (arithmetic) as appropriate for businessmen and men of war who "must learn the art of numbers or [they] will not know how to array [their] troops" and arithmetic (number theory) as appropriate for philosophers "because [they have] to arise out of the sea of change and lay hold of true being." Euclid of Alexandria, when asked by one of his students of what use was the study of geometry, asked his slave to give the student threepence, "since he must make gain of what he learns." The Greek mathematician Apollonius of Perga was asked about the usefulness of some of his theorems in Book IV of Conics to which he proudly asserted,
They are worthy of acceptance for the sake of the demonstrations themselves, in the same way as we accept many other things in mathematics for this and for no other reason.
And since many of his results were not applicable to the science or engineering of his day, Apollonius further argued in the preface of the fifth book of Conics that the subject is one of those that "...seem worthy of study for their own sake."
19th century
The term itself is enshrined in the full title of the Sadleirian Chair, "Sadleirian Professor of Pure Mathematics", founded (as a professorship) in the mid-nineteenth century. The idea of a separate discipline of pure mathematics may have emerged at that time. The generation of Gauss made no sweeping distinction of the kind between pure and applied. In the following years, specialisation and professionalisation (particularly in the Weierstrass approach to mathematical analysis) started to make a rift more apparent.
20th century
At the start of the twentieth century mathematicians took up the axiomatic method, strongly influenced by David Hilbert's example. The logical formulation of pure mathematics suggested by Bertrand Russell in terms of a quantifier structure of propositions seemed more and more plausible, as large parts of mathematics became axiomatised and thus subject to the simple criteria of rigorous proof.
Pure mathematics, according to a view that can be ascribed to the Bourbaki group, is what is proved. "Pure mathematician" became a recognized vocation, achievable through training.
The case was made that pure mathematics is useful in engineering education:
There is a training in habits of thought, points of view, and intellectual comprehension of ordinary engineering problems, which only the study of higher mathematics can give.
Generality and abstraction
One central concept in pure mathematics is the idea of generality; pure mathematics often exhibits a trend towards increased generality. Uses and advantages of generality include the following:
Generalizing theorems or mathematical structures can lead to deeper understanding of the original theorems or structures
Generality can simplify the presentation of material, resulting in shorter proofs or arguments that are easier to follow.
One can use generality to avoid duplication of effort, proving a general result instead of having to prove separate cases independently, or using results from other areas of mathematics.
Generality can facilitate connections between different branches of mathematics. Category theory is one area of mathematics dedicated to exploring this commonality of structure as it plays out in some areas of math.
Generality's impact on intuition is both dependent on the subject and a matter of personal preference or learning style. Often generality is seen as a hindrance to intuition, although it can certainly function as an aid to it, especially when it provides analogies to material for which one already has good intuition.
As a prime example of generality, the Erlangen program involved an expansion of geometry to accommodate non-Euclidean geometries as well as the field of topology, and other forms of geometry, by viewing geometry as the study of a space together with a group of transformations. The study of numbers, called algebra at the beginning undergraduate level, extends to abstract algebra at a more advanced level; and the study of functions, called calculus at the college freshman level becomes mathematical analysis and functional analysis at a more advanced level. Each of these branches of more abstract mathematics have many sub-specialties, and there are in fact many connections between pure mathematics and applied mathematics disciplines. A steep rise in abstraction was seen mid 20th century.
In practice, however, these developments led to a sharp divergence from physics, particularly from 1950 to 1983. Later this was criticised, for example by Vladimir Arnold, as too much Hilbert, not enough Poincaré. The point does not yet seem to be settled, in that string theory pulls one way, while discrete mathematics pulls back towards proof as central.
Pure vs. applied mathematics
Mathematicians have always had differing opinions regarding the distinction between pure and applied mathematics. One of the most famous (but perhaps misunderstood) modern examples of this debate can be found in G.H. Hardy's 1940 essay A Mathematician's Apology.
It is widely believed that Hardy considered applied mathematics to be ugly and dull. Although it is true that Hardy preferred pure mathematics, which he often compared to painting and poetry, Hardy saw the distinction between pure and applied mathematics to be simply that applied mathematics sought to express physical truth in a mathematical framework, whereas pure mathematics expressed truths that were independent of the physical world. Hardy made a separate distinction in mathematics between what he called "real" mathematics, "which has permanent aesthetic value", and "the dull and elementary parts of mathematics" that have practical use.
Hardy considered some physicists, such as Einstein and Dirac, to be among the "real" mathematicians, but at the time that he was writing his Apology, he considered general relativity and quantum mechanics to be "useless", which allowed him to hold the opinion that only "dull" mathematics was useful. Moreover, Hardy briefly admitted that—just as the application of matrix theory and group theory to physics had come unexpectedly—the time may come where some kinds of beautiful, "real" mathematics may be useful as well.
Another insightful view is offered by American mathematician Andy Magid:
Friedrich Engels argued in his 1878 book Anti-Dühring that "it is not at all true that in pure mathematics the mind deals only with its own creations and imaginations. The concepts of number and figure have not been invented from any source other than the world of reality". He further argued that "Before one came upon the idea of deducing the form of a cylinder from the rotation of a rectangle about one of its sides, a number of real rectangles and cylinders, however imperfect in form, must have been examined. Like all other sciences, mathematics arose out of the needs of men...But, as in every department of thought, at a certain stage of development the laws, which were abstracted from the real world, become divorced from the real world, and are set up against it as something independent, as laws coming from outside, to which the world has to conform."
| Mathematics | Other | null |
457601 | https://en.wikipedia.org/wiki/Salmonella%20enterica | Salmonella enterica | Salmonella enterica (formerly Salmonella choleraesuis) is a rod-shaped, flagellate, facultative anaerobic, Gram-negative bacterium and a species of the genus Salmonella. It is divided into six subspecies, arizonae (IIIa), diarizonae (IIIb), houtenae (IV), salamae (II), indica (VI), and enterica (I). A number of its serovars are serious human pathogens; many of them are (more specifically) serovars of Salmonella enterica subsp. enterica.
Epidemiology
Most cases of salmonellosis are caused by food infected with S. enterica, which often infects cattle and poultry, though other animals such as domestic cats and hamsters have also been shown to be sources of infection in humans. It primarily resides in the intestinal tract of animals and humans and can be found in feedstuff, soil, bedding, litter, and fecal matter.
The primary reservoir for the pathogen is poultry and 70% of human cases are attributed to the consumption of contaminated eggs, chicken, or turkey. Raw chicken eggs and goose eggs can harbor S. enterica, initially in the egg whites, although most eggs are not infected. As the egg ages at room temperature, the yolk membrane begins to break down and S. enterica can spread into the yolk. Refrigeration and freezing do not kill all the bacteria, but substantially slow or halt their growth. Pasteurizing and food irradiation are used to kill Salmonella for commercially produced foodstuffs containing raw eggs such as ice cream. Foods prepared in the home from raw eggs, such as mayonnaise, cakes, and cookies, can spread salmonellae if not properly cooked before consumption. Salmonella is the leading foodborne pathogen in the United States, causing the most deaths and having the highest cost burden. It is a resilient microorganism capable of surviving long periods of time in hot and dry environments, increasing its effectiveness as a pathogen and making it able to survive the harsh environments of the gastrointestinal tract and farms. Salmonella has been found in 10% to 26% of farm environments in Tennessee, North Carolina, Alabama, California, and Washington.
S. enterica genomes have been reconstructed from up to 6,500 year old human remains across Western Eurasia, which provides evidence for geographic widespread infections with systemic S. enterica during prehistory, and a possible role of the Neolithization process in the evolution of host adaptation. Additional reconstructed genomes from colonial Mexico suggest S. enterica as the cause of cocoliztli, an epidemic in 16th-century New Spain. In 1545, this outbreak of S. enterica spread explosively across what is now Mexico. Over the next century, the disease killed up to 90% of the Indigenous population.
Children under the age of five years, the elderly, and immunosuppressed adults are at an increased risk of systemic dissemination of the disease and need specialized treatment to combat the disease. Drinking extra fluids and antibiotics such as fluoroquinolones are typical treatments. Complications of the disease often appear as anemia or septicaemia, and the mortality rate is 15% once these symptoms arise.
The serogroup S. Typhi is the cause of typhoid fever.
Nomenclature
S. enterica has six subspecies, and each subspecies has associated serovars that differ by antigenic specificity. S. enterica has over 2500 serovars. Salmonella bongori was previously considered a subspecies of S. enterica, but it is now the other species in the genus Salmonella. Most of the human pathogenic Salmonella serovars belong to the enterica subspecies. These serogroups include S. Typhi, S. Enteritidis, S. Paratyphi, S. Typhimurium, and S. Choleraesuis. The serovars can be designated as written in the previous sentence (capitalized and nonitalicized following the genus), or as follows: "S. enterica subsp. enterica, serovar Typhi".
Subspecies S. e. arizonae, named after the state of Arizona, is most commonly found in cold-blooded animals (especially snakes), but can also infect turkey, sheep, and humans. It is endemic in southwestern United States. The similar S. e. subsp. diarizonae also infects snakes and occasionally humans.
Pathogenesis
Secreted proteins are of major importance for the pathogenesis of infectious diseases caused by S. enterica. A remarkably large number of fimbrial and nonfimbrial adhesins are present in Salmonella, and mediate biofilm formation and contact to host cells. Secreted proteins are also involved in host-cell invasion and intracellular proliferation, two hallmarks of Salmonella pathogenesis.
DNA repair capability
Exposure of S. enterica to bile salts, such as sodium deoxycholate, induces the SOS DNA damage response indicating that in this organism bile salts cause DNA damage. Bile salt exposure is found to increase GC to AT transition mutations and also to induce genes of the OxyR and SoxRS regulons suggesting further that bile salts specifically cause oxidative DNA damage. Mutants of S. enterica that are defective in enzymes required for the process of base excision repair are sensitive to bile salts. This indicates that wild-type S. enterica uses base excision repair to remove DNA damages caused by the bile salts. The RecBCD enzyme which functions in recombinational repair of DNA is also required for bile salt resistance.
Small noncoding RNA
Small nonprotein-coding RNAs (sRNA) are able to perform specific functions without being translated into proteins; 97 bacterial sRNAs from Salmonella Typhi were discovered.
AsdA (antisense RNA of dnaA) is a cis-encoded antisense RNA of dnaA described in S. enterica serovar Typhi. It was discovered by deep sequencing and its transcription was confirmed by Northern blot and RACE analysis. AsdA is estimated to be about 540 nucleotides long, and represents the complementary strand to that encoding DnaA, a protein that plays a central role in the initiation of DNA replication and hence cellular division. In rich media, it is highly expressed only after reaching the stationary growth phase, but under limiting iron or osmotic stress, it is already expressed during exponential growth. Overexpression of AsdA stabilizes dnaA mRNA, increasing its levels and thereby enhancing its rate of translation. This suggests that AsdA is a regulator of DNA replication.
| Biology and health sciences | Gram-negative bacteria | Plants |
457857 | https://en.wikipedia.org/wiki/Traditional%20medicine | Traditional medicine | Traditional medicine (also known as indigenous medicine or folk medicine) comprises medical aspects of traditional knowledge that developed over generations within the folk beliefs of various societies, including indigenous peoples, before the era of modern medicine. The World Health Organization (WHO) defines traditional medicine as "the sum total of the knowledge, skills, and practices based on the theories, beliefs, and experiences indigenous to different cultures, whether explicable or not, used in the maintenance of health as well as in the prevention, diagnosis, improvement and treatment of physical and mental illness". Traditional medicine is often contrasted with Evidence based medicine.
In some Asian and African countries, up to 80% of the population relies on traditional medicine for their primary health care needs. Traditional medicine is a form of alternative medicine. Practices known as traditional medicines include traditional European medicine, traditional Chinese medicine, traditional Korean medicine, traditional African medicine, Ayurveda, Siddha medicine, Unani, ancient Iranian medicine, traditional Iranian medicine, medieval Islamic medicine, Muti, Ifá and Rongoā. Scientific disciplines that study traditional medicine include herbalism, ethnomedicine, ethnobotany, and medical anthropology.
The WHO notes, however, that "inappropriate use of traditional medicines or practices can have negative or dangerous effects" and that "further research is needed to ascertain the efficacy and safety" of such practices and medicinal plants used by traditional medicine systems. Its "Traditional Medicine Strategy 2014–2023" said that the WHO would "support Member States in developing proactive policies and implementing action plans that will strengthen the role traditional medicine plays in keeping populations healthy."
Usage and history
Classical history
In the written record, the study of herbs dates back 5,000 years to the ancient Sumerians, who described well-established medicinal uses for plants. In Ancient Egyptian medicine, the Ebers papyrus from c. 1552 BC records a list of folk remedies and magical medical practices. The Old Testament also mentions herb use and cultivation in regards to Kashrut.
Many herbs and minerals used in Ayurveda were described by ancient Indian herbalists such as Charaka and Sushruta during the 1st millennium BC. The first Chinese herbal book was the Shennong Bencaojing, compiled during the Han dynasty but dating back to a much earlier date, which was later augmented as the Yaoxing Lun (Treatise on the Nature of Medicinal Herbs) during the Tang dynasty. Early recognised Greek compilers of existing and current herbal knowledge include Pythagoras and his followers, Hippocrates, Aristotle, Theophrastus, Dioscorides and Galen.
Roman sources included Pliny the Elder's Natural History and Celsus's De Medicina. Pedanius Dioscorides drew on and corrected earlier authors for his De Materia Medica, adding much new material; the work was translated into several languages, and Turkish, Arabic and Hebrew names were added to it over the centuries. Latin manuscripts of De Materia Medica were combined with a Latin herbal by Apuleius Platonicus (Herbarium Apuleii Platonici) and were incorporated into the Anglo-Saxon codex Cotton Vitellius C.III. These early Greek and Roman compilations became the backbone of European medical theory and were translated by the Persian Avicenna (Ibn Sīnā, 980–1037), the Persian Rhazes (Rāzi, 865–925) and the Jewish Maimonides.
Some fossils have been used in traditional medicine since antiquity.
Medieval and later
Arabic indigenous medicine developed from the conflict between the magic-based medicine of the Bedouins and the Arabic translations of the Hellenic and Ayurvedic medical traditions. Spanish medicine was influenced by the Arabs from 711 to 1492. Islamic physicians and Muslim botanists such as al-Dinawari and Ibn al-Baitar significantly expanded on the earlier knowledge of materia medica. The most famous Persian medical treatise was Avicenna's The Canon of Medicine, which was an early pharmacopoeia and introduced clinical trials. The Canon was translated into Latin in the 12th century and remained a medical authority in Europe until the 17th century. The Unani system of traditional medicine is also based on the Canon.
Translations of the early Roman-Greek compilations were made into German by Hieronymus Bock whose herbal, published in 1546, was called Kreuter Buch. The book was translated into Dutch as Pemptades by Rembert Dodoens (1517–1585), and from Dutch into English by Carolus Clusius, (1526–1609), published by Henry Lyte in 1578 as A Nievve Herball. This became John Gerard's (1545–1612) Herball or General Historie of Plantes. Each new work was a compilation of existing texts with new additions.
Women's folk knowledge existed in undocumented parallel with these texts. Forty-four drugs, diluents, flavouring agents and emollients mentioned by Dioscorides are still listed in the official pharmacopoeias of Europe. The Puritans took Gerard's work to the United States where it influenced American Indigenous medicine.
Francisco Hernández, physician to Philip II of Spain spent the years 1571–1577 gathering information in Mexico and then wrote Rerum Medicarum Novae Hispaniae Thesaurus, many versions of which have been published including one by Francisco Ximénez. Both Hernandez and Ximenez fitted Aztec ethnomedicinal information into the European concepts of disease such as "warm", "cold", and "moist", but it is not clear that the Aztecs used these categories. Juan de Esteyneffer's Florilegio medicinal de todas las enfermedas compiled European texts and added 35 Mexican plants.
Martín de la Cruz wrote a herbal in Nahuatl which was translated into Latin by Juan Badiano as Libellus de Medicinalibus Indorum Herbis or Codex Barberini, Latin 241 and given to King Carlos V of Spain in 1552. It was apparently written in haste and influenced by the European occupation of the previous 30 years. Fray Bernardino de Sahagún's used ethnographic methods to compile his codices that then became the Historia General de las Cosas de Nueva España, published in 1793. Castore Durante published his Herbario Nuovo in 1585 describing medicinal plants from Europe and the East and West Indies. It was translated into German in 1609 and Italian editions were published for the next century.
Colonial America
In 17th and 18th-century America, traditional folk healers, frequently women, used herbal remedies, cupping and leeching. Native American traditional herbal medicine introduced cures for malaria, dysentery, scurvy, non-venereal syphilis, and goiter problems. Many of these herbal and folk remedies continued on through the 19th and into the 20th century, with some plant medicines forming the basis for modern pharmacology.
Modern usage
The prevalence of folk medicine in certain areas of the world varies according to cultural norms. Some modern medicine is based on plant phytochemicals that had been used in folk medicine. Researchers state that many of the alternative treatments are "statistically indistinguishable from placebo treatments".
Knowledge transmission and creation
Indigenous medicine is generally transmitted orally through a community, family and individuals until "collected". Within a given culture, elements of indigenous medicine knowledge may be diffusely known by many, or may be gathered and applied by those in a specific role of healer such as a shaman or midwife. Three factors legitimize the role of the healer – their own beliefs, the success of their actions and the beliefs of the community. When the claims of indigenous medicine become rejected by a culture, generally three types of adherents still use it – those born and socialized in it who become permanent believers, temporary believers who turn to it in crisis times, and those who only believe in specific aspects, not in all of it.
Definition and terminology
Traditional medicine may sometimes be considered as distinct from folk medicine, and considered to include formalized aspects of folk medicine. Under this definition folk medicine are longstanding remedies and practises passed on and practiced by lay people. Folk medicine consists of the healing modalities, ideas of body physiology and health preservation known to some in a culture, transmitted informally as general knowledge, and practiced or applied by anyone in the culture having prior experience.
Folk medicine
Many countries have practices described as folk medicine which may coexist with formalized, science-based, and institutionalized systems of medical practice represented by conventional medicine. Examples of folk medicine traditions are traditional Chinese medicine, Iranian traditional medicine, traditional Korean medicine, Arabic indigenous medicine, Uyghur traditional medicine, Japanese Kampō medicine, traditional Aboriginal bush medicine, Native Hawaiian Lāʻau lapaʻau, Curanderismo norteño, and Georgian folk medicine, among others.
Australian bush medicine
Generally, bush medicine used by Aboriginal and Torres Strait Islander people in Australia is made from plant materials, such as bark, leaves and seeds, although animal products may be used as well. A major component of traditional medicine is herbal medicine, which is the use of natural plant substances to treat or prevent illness.
Native American medicine
American Native and Alaska Native medicine are traditional forms of healing that have been around for thousands of years. There are many ethnobotany plants involved in traditional medicine for Native Americans and some are still used today. When it comes to Native American traditional medicine, the ideas surrounding health and illness within the culture are virtually inseparable from the ideas of religion and spirituality. Healers within indigenous communities go by many names ranging from medicine man or woman to herbalist or even shaman and are considered spiritual or religious leaders within their respective tribes. When it comes to healing, tribal healers would look at a plant's characteristics to determine its efficacy for the treatment of an illness. Specific plant characteristics such as plant shape, smell, color, and taste could aid in determining how the plant could be used as a remedy. The Meskwaki tribe found they could use the juice from Arum maculatum for snakebites. This was inferred from the milky appearance of the juice from the plant which is said to resemble snake venom, and the plant's shape resembled the head of a snake. Native Americans used foxglove herb as a treatment for an illness they referred to as dropsy or edema, which is fluid buildup typically in the lower legs, and its common cause is heart failure. In modern medicine, foxglove extract is still used under the name digitalis, and its purpose is to moderate the heart rate. Native Americans were successful with some medical practices, such as treating fevers, gastrointestinal conditions, skin rashes, setting bones, as well as birthing babies, and aiding mothers in healing. A study conducted within an IHS hospital that allows Navajo healers to visit patients found that the hospital had an 80 percent success rate in getting comatose patients back to consciousness, which is higher than the rate of present-day biomedical management hospitals. The plant family Asteraceae has been commonly selected for orthopedic aids and pulmonary aids, specifically the species Achillea and Artemisia. A study conducted amongst 14 different tribes within North America found that Asteraceae was the most widely used plant family for its medicinal properties.
Nattuvaidyam
Nattuvaidyam was a set of indigenous medical practices that existed in India before the advent of allopathic or western medicine. These practices had different sets of principles and ideas of the body, health and disease. There were overlaps and borrowing of ideas, medicinal compounds used and techniques within these practices. Some of these practices had written texts in vernacular languages like Malayalam, Tamil, Telugu, etc. while others were handed down orally through various mnemonic devices. Ayurveda was one kind of nattuvaidyam practised in south India. The others were kalarichikitsa (related to bone setting and musculature), marmachikitsa (vital spot massaging), ottamoolivaidyam (single dose medicine or single time medication), chintamanivaidyam and so on. When the medical system was revamped in twentieth century India, many of the practices and techniques specific to some of these diverse nattuvaidyam were included in Ayurveda.
Home remedies
A home remedy (sometimes also referred to as a granny cure) is a treatment to cure a disease or ailment that employs certain spices, herbs, vegetables, or other common items. Home remedies may or may not have medicinal properties that treat or cure the disease or ailment in question, as they are typically passed along by laypersons (which has been facilitated in recent years by the Internet). Many are merely used as a result of tradition or habit or because they are effective in inducing the placebo effect.
One of the more popular examples of a home remedy is the use of chicken soup as an aid in treating respiratory infections such as a cold or mild flu. Other examples of home remedies include duct tape to help with setting broken bones; duct tape or superglue to treat plantar warts; and Kogel mogel to treat sore throat. In earlier times, mothers were entrusted with all but serious remedies. Historic cookbooks are frequently full of remedies for dyspepsia, fevers, and female complaints. Components of the aloe vera plant are used to treat skin disorders. Many European liqueurs or digestifs were originally sold as medicinal remedies. In Chinese folk medicine, medicinal congees (long-cooked rice soups with herbs), foods, and soups are part of treatment practices.
Criticism
Safety concerns
Although 130 countries have regulations on folk medicines, there are risks associated with the use of them (i.e. zoonosis, mainly as some traditional medicines still use animal-based substances). It is often assumed that because supposed medicines are natural that they are safe, but numerous precautions are associated with using herbal remedies.
Use of endangered species
Endangered animals, such as the slow loris, are sometimes killed to make traditional medicines.
Shark fins have also been used in traditional medicine, and although their effectiveness has not been proven, it is hurting shark populations and their ecosystem.
The illegal ivory trade can partially be traced back to buyers of traditional Chinese medicine. Demand for ivory is a huge factor in the poaching of endangered species such as rhinos and elephants.
Pangolins are threatened by poaching for their meat and scales, which are used in traditional medicine. They are the most trafficked mammals in the world.
| Biology and health sciences | Alternative and traditional medicine | null |
457921 | https://en.wikipedia.org/wiki/Fringe%20science | Fringe science | Fringe science refers to ideas whose attributes include being highly speculative or relying on premises already refuted. Fringe science theories are often advanced by people who have no traditional academic science background, or by researchers outside the mainstream discipline. The general public has difficulty distinguishing between science and its imitators, and in some cases, a "yearning to believe or a generalized suspicion of experts is a very potent incentive to accepting pseudoscientific claims".
The term "fringe science" covers everything from novel hypotheses, which can be tested utilizing the scientific method, to wild ad hoc hypotheses and mumbo jumbo. This has resulted in a tendency to dismiss all fringe science as the domain of pseudoscientists, hobbyists, and quacks.
A concept that was once accepted by the mainstream scientific community may become fringe science because of a later evaluation of previous research. For example, focal infection theory, which held that focal infections of the tonsils or teeth are a primary cause of systemic disease, was once considered to be medical fact. It has since been dismissed because of a lack of evidence.
Description
The boundary between fringe science and pseudoscience is disputed. The connotation of "fringe science" is that the enterprise is rational but is unlikely to produce good results for various reasons, including incomplete or contradictory evidence. Pseudoscience, however, is something that is not scientific but is incorrectly characterised as science.
The term may be considered pejorative. For example, Lyell D. Henry Jr. wrote, "Fringe science [is] a term also suggesting kookiness." This characterization is perhaps inspired by the eccentric behavior of many researchers of the kind known colloquially (and with considerable historical precedent) as mad scientists.
Although most fringe science is rejected, the scientific community has come to accept some portions of it. One example of such is plate tectonics, an idea which had its origin in the fringe science of continental drift and was rejected for decades.
Examples
Historical
Some historical ideas that are considered to have been refuted by mainstream science are:
Wilhelm Reich's work with orgone, a physical energy he claimed to have discovered, contributed to his alienation from the psychiatric community. He was eventually sentenced to two years in a federal prison, where he died. At that time and continuing today, scientists disputed his claim that he had scientific evidence for the existence of orgone. Nevertheless, amateurs and a few fringe researchers continued to believe that orgone is real.
Focal infection theory (FIT), as the primary cause of systemic disease, rapidly became accepted by mainstream dentistry and medicine after World War I. This acceptance was largely based upon what later turned out to be fundamentally flawed studies. As a result, millions of people were subjected to needless dental extractions and surgeries. The original studies supporting FIT began falling out of favor in the 1930s. By the late 1950s, it was regarded as a fringe theory.
The Clovis First theory held that the Clovis culture was the first culture in North America. It was long regarded as a mainstream theory until mounting evidence of a pre-Clovis culture discredited it.
Modern
Relatively recent fringe sciences include:
Aubrey de Grey, featured in a 2006 60 Minutes special report, is studying human longevity. He calls his work "strategies for engineered negligible senescence" (SENS). Many mainstream scientists believe his research is fringe science (especially his view of the importance of nuclear epimutations and his timeline for antiaging therapeutics). In a 2005 article in Technology Review (part of a larger series), it was stated that "SENS is highly speculative. Many of its proposals have not been reproduced, nor could they be reproduced with today's scientific knowledge and technology. Echoing Myhrvold, we might charitably say that de Grey's proposals exist in a kind of antechamber of science, where they wait (possibly in vain) for independent verification. SENS does not compel the assent of many knowledgeable scientists; but neither is it demonstrably wrong."
A nuclear fusion reaction called cold fusion, which occurs near room temperature and pressure, was reported by chemists Martin Fleischmann and Stanley Pons in March 1989. Numerous research efforts at the time were unable to replicate their results. Subsequently, several scientists have worked on cold fusion or have participated in international conferences on it. In 2004, the United States Department of Energy commissioned a panel on cold fusion to reexamine the concept. They wanted to determine whether their policies should be altered because of new evidence.
The theory of abiogenic petroleum origin holds that petroleum was formed from deep carbon deposits, perhaps dating to the formation of the Earth. The ubiquity of hydrocarbons in the solar system may be evidence that there may be more petroleum on Earth than commonly thought and that petroleum may originate from carbon-bearing fluids that migrate upward from the Earth's mantle. Abiogenic hypotheses saw a revival in the last half of the twentieth century by Russian and Ukrainian scientists. More interest was generated in the West after the 1999 publication by Thomas Gold of The Deep Hot Biosphere. Gold's version of the theory is partly based on the existence of a biosphere composed of thermophile bacteria in the Earth's crust, which might explain the existence of specific biomarkers in extracted petroleum.
Accepted as mainstream
Some theories that were once rejected as fringe science but were eventually accepted as mainstream science include:
Plate tectonics
The existence of Troy
Heliocentrism
Norse colonization of the Americas
The Big Bang theory
Helicobacter pylori bacteria as the causative agent of peptic ulcer disease
The germ theory of disease
Neanderthal-Homo sapiens hybridization
Responding to fringe science
Michael W. Friedlander has suggested some guidelines for responding to fringe science, which, he argues, is a more difficult problem than scientific misconduct. His suggested methods include impeccable accuracy, checking cited sources, not overstating orthodox science, thorough understanding of the Wegener continental drift example, examples of orthodox science investigating radical proposals, and prepared examples of errors from fringe scientists.
Friedlander suggests that fringe science is necessary so mainstream science will not atrophy. Scientists must evaluate the plausibility of each new fringe claim, and certain fringe discoveries "will later graduate into the ranks of accepted" — while others "will never receive confirmation".
Margaret Wertheim profiled many "outsider scientists" in her book Physics on the Fringe, who receive little or no attention from professional scientists. She describes all of them as trying to make sense of the world using the scientific method but in the face of being unable to understand modern science's complex theories. She also finds it fair that credentialed scientists do not bother spending a lot of time learning about and explaining problems with the fringe theories of uncredentialed scientists since the authors of those theories have not taken the time to understand the mainstream theories they aim to disprove.
Controversies
As Donald E. Simanek asserts, "Too often speculative and tentative hypotheses of cutting edge science are treated as if they were scientific truths, and so accepted by a public eager for answers." However, the public is ignorant that "As science progresses from ignorance to understanding it must pass through a transitional phase of confusion and uncertainty."
The media also play a role in propagating the belief that certain fields of science are controversial. In their 2003 paper "Optimising Public Understanding of Science and Technology in Europe: A Comparative Perspective", Jan Nolin et al. write that "From a media perspective it is evident that controversial science sells, not only because of its dramatic value, but also since it is often connected to high-stake societal issues."
| Physical sciences | Science basics | Basics and measurement |
457926 | https://en.wikipedia.org/wiki/B%20vitamins | B vitamins | B vitamins are a class of water-soluble vitamins that play important roles in cell metabolism and synthesis of red blood cells. They are a chemically diverse class of compounds.
Dietary supplements containing all eight are referred to as a vitamin B complex. Individual B vitamins are referred to by B-number or by chemical name, such as B1 for thiamine, B2 for riboflavin, and B3 for niacin, while some are more commonly recognized by name than by number, such as pantothenic acid (B5), biotin (B7), and folate (B9). B vitamins are present in protein-rich foods, such as fish, poultry, meat, dairy products, and eggs; they are also found in leafy green vegetables, beans, and peas. Fortified foods, such as breakfast cereals, baked products, and infant formulas, may contain B vitamins.
Each B vitamin is either a cofactor (generally a coenzyme) for key metabolic processes or is a precursor needed to make one.
List of B vitamins
Note: Other substances once thought to be vitamins were given B-numbers, but were disqualified once discovered to be either manufactured by the body or not essential for life. See #Related compounds for numbers 4, 8, 10, 11, and others.
Sources
B vitamins are found in abundance in meat, eggs, and dairy products. Processed carbohydrates such as sugar and white flour tend to have lower B vitamin content than their unprocessed counterparts. For this reason, it is common in many countries (including the United States) that the B vitamins thiamine, riboflavin, niacin, and folic acid are added back to white flour after processing. This is referred to as "enriched flour" on food labels. B vitamins are particularly concentrated in meat such as turkey, tuna and liver.
Sources for B vitamins also include spinach, legumes (pulses or beans), whole grains, asparagus, potatoes, bananas, chili peppers, breakfast cereals.
The B12 vitamin is not abundantly available from plant products (although it has been found in moderate abundance in fermented vegetable products, certain seaweeds, and in certain mushrooms, with the bioavailability of the vitamin in these cases remaining uncertain), making B12 deficiency a legitimate concern for those maintaining a vegan diet. Manufacturers of plant-based foods will sometimes report B12 content, leading to confusion about what sources yield B12. The confusion arises because the standard US Pharmacopeia (USP) method for measuring the B12 content does not measure the B12 directly. Instead, it measures a bacterial response to the food. Chemical variants of the B12 vitamin found in plant sources are active for bacteria, but cannot be used by the human body. This same phenomenon can cause significant over-reporting of B12 content in other types of foods as well.
A common way to increase vitamin B intake is by using dietary supplements. B vitamins are commonly added to energy drinks, many of which have been marketed with large amounts of B vitamins.
Because they are soluble in water, excess B vitamins are generally readily excreted, although individual absorption, use and metabolism may vary. The elderly and athletes may need to supplement their intake of B12 and other B vitamins due to problems in absorption and increased needs for energy production. In cases of severe deficiency, B vitamins, especially B12, may also be delivered by injection to reverse deficiencies. Both type 1 and type 2 diabetics may also be advised to supplement thiamine based on high prevalence of low plasma thiamine concentration and increased thiamine clearance associated with diabetes. Also, folate deficiency in early embryo development has been linked to neural tube defects. Thus, women planning to become pregnant are usually encouraged to increase daily dietary folate intake or take a supplement.
Molecular functions
To the right, a diagram of some of the major B vitamins (2, 3, 5, 9, and 12) are shown as precursors for certain essential biochemical reactants (FAD, NAD+, coenzyme A, and heme B respectively). The structural similarities between them are highlighted, which illustrates the precursor nature of many B vitamins while also showing the functionality of the end product used by essential reactions to support human, animal, or cellular life.
FAD, NAD+, and coenzyme A are all essential for the catabolic release of free energy (dG) to power the activity of the cell and more complex life forms. See the article on Catabolism for more details on how these three essential biochemical reactants help support life.
Tetrahydrofolate is a necessary co-reactant for synthesizing some amino acids, such as glycine. Heme B is the porphyrin derivative macrocycle molecule that holds the iron atom in place in hemoglobin, allowing for the transportation of oxygen through blood.
Deficiencies
Several named vitamin deficiency diseases may result from the lack of sufficient B vitamins. Deficiencies of other B vitamins result in symptoms that are not part of a named deficiency disease.
Side effects
Because water-soluble B vitamins are eliminated in the urine, taking large doses of certain B vitamins usually only produces transient side effects (only exception is pyridoxine). General side effects may include restlessness, nausea and insomnia. These side effects are almost always caused by dietary supplements and not foodstuffs.
Discovery
Related compounds
Many of the following substances have been referred to as vitamins as they were once believed to be vitamins. They are no longer considered as such, and the numbers that were assigned to them now form the "gaps" in the true series of B-complex vitamins described above (for example, there is no vitamin B4). Some of them, though not essential to humans, are essential in the diets of other organisms; others have no known nutritional value and may even be toxic under certain conditions.
Vitamin B4: can refer to the distinct chemicals choline, adenine, or carnitine.
Choline is synthesized by the human body, but not sufficiently to maintain good health, and is now considered an essential dietary nutrient.
Adenine is a nucleobase synthesized by the human body.
Carnitine is an essential dietary nutrient for certain worms, but not for humans.
Vitamin B8: adenosine monophosphate (AMP), also known as adenylic acid. Vitamin B8 may also refer to inositol.
Vitamin B10: para-aminobenzoic acid (pABA or PABA), a chemical component of the folate molecule produced by plants and bacteria, and found in many foods. It is best known as a UV-blocking sunscreen applied to the skin, and is sometimes taken orally for certain medical conditions.
Vitamin B11: pteroylheptaglutamic acid (PHGA; chick growth factor). Vitamin Bc-conjugate was also found to be identical to PHGA. Derivative of folate ("pteroylmonoglutamic acid" in this nomenclature).
Vitamin B13: orotic acid.
Vitamin B14: cell proliferant, anti-anemia, rat growth factor, and antitumor pterin phosphate, named by Earl R. Norris. Isolated from human urine at 0.33ppm (later in blood), but later abandoned by him as further evidence did not confirm this. He also claimed this was not xanthopterin.
Vitamin B15: pangamic acid, also known as pangamate. Promoted in various forms as a dietary supplement and drug; considered unsafe and subject to seizure by the US Food and Drug Administration.
Vitamin B16: dimethylglycine (DMG) is synthesized by the human body from choline.
Vitamin B17: pseudoscientific name for the poisonous compound amygdalin, also known as the equally pseudoscientific name "nitrilosides" despite the fact that it is a single compound. Amygdalin can be found in various plants, but is most commonly extracted from apricot pits and other similar fruit kernels. Amygdalin is hydrolyzed by various intestinal enzymes to form, among other things, hydrogen cyanide, which is toxic to human beings when exposed to a high enough dosage. Some proponents claim that amygdalin is effective in cancer treatment and prevention, despite its toxicity and a lack of scientific evidence.
Vitamin B20: L-carnitine.
Vitamin Bf: carnitine.
Vitamin Bm: myo-inositol, also called "mouse antialopaecia factor".
Vitamin Bp: "antiperosis factor", which prevents perosis, a leg disorder, in chicks; can be replaced by choline and manganese salts.
Vitamin BT: carnitine.
Vitamin Bv: a type of B6 other than pyridoxine.
Vitamin BW: a type of biotin other than d-biotin.
Vitamin Bx: an alternative name for both pABA (see vitamin B10) and pantothenic acid.
| Biology and health sciences | Vitamins | Health |
457970 | https://en.wikipedia.org/wiki/Podocarpaceae | Podocarpaceae | Podocarpaceae is a large family of mainly Southern Hemisphere conifers, known in English as podocarps, comprising about 156 species of evergreen trees and shrubs. It contains 19 genera if Phyllocladus is included and Manoao and Sundacarpus are recognized. The family achieved its maximum diversity in the Cenozoic, making the Podocarpaceae family one of the most diverse in the southern hemisphere.
The family is a classic member of the Antarctic flora, with its main centres of diversity in Australasia, particularly New Caledonia, Tasmania, and New Zealand, and to a slightly lesser extent Malesia and South America (primarily in the Andes Mountains). Several genera extend north of the equator into Indochina and the Philippines. Podocarpus reaches as far north as southern Japan and southern China in Asia, and Mexico in the Americas, and Nageia into southern China and southern India. Two genera also occur in sub-Saharan Africa, the widespread Podocarpus and the endemic Afrocarpus.
Parasitaxus usta is unique as the only known parasitic gymnosperm. It occurs on New Caledonia, where it is parasitic on another member of the Podocarpaceae, Falcatifolium taxoides.
The genus Phyllocladus is sister to the Podocarpaceae sensu stricto. It is treated by some botanists in its own family, the Phyllocladaceae.
Taxonomy
The Podocarpaceae show great diversity, both morphologically and ecologically. Members occur mainly in the Southern Hemisphere, with most genetic variety taking place in New Caledonia, New Zealand, and Tasmania. Species diversity of Podocarpus is found mainly in South America and the Indonesian islands, the latter also being rich in Dacrydium and Dacrycarpus species.
Podocarpus (with 82 to 100 species) and Dacrydium (with 21 species) are the largest genera. A few genera are common to New Zealand and South America, supporting the view that podocarps had an extensive distribution over southern Gondwanaland. The breaking up of Gondwanaland led to large-scale speciation of the Podocarpaceae.
Until 1970, only seven Podocarpaceae genera were recognized: Podocarpus, Dacrydium, Phyllocladus, Acmopyle, Microcachrys, Saxegothaea, and Pherosphaera. All four of the African species fell under Podocarpus – P. falcatus, P. elongatus, P. henkelii, and P. latifolius. Taxonomists divided Podocarpus species into eight species groups based on leaf anatomy: Afrocarpus J.Buchholz & N.E.Gray, Dacrycarpus Endl., Eupodocarpus Endl., Microcarpus Pilg., Nageia (Gaertn.) Endl., Polypodiopsis C.E.Bertrand (non Polypodiopsis Carriére nom. rej. prop. 6), Stachycarpus Endl. and Sundacarpus J.Buchholz and N.E.Gray.
Studies of embryology, gametophyte development, female cone structure, and cytology led to the belief that the eight categories probably deserved generic status. Researchers agreed on the need to recognize "fairly natural groupings which prove to have good geographic and probably evolutionary cohesion" and took the necessary steps to raise each section to generic status.
In 1990, a treatment of the Podocarpaceae recognized 17 genera, excluding Phyllocladus from the family, while recognizing Sundacarpus, but not Manoao. In 1995, Manoao was segregated from Lagarostrobus, based on morphological characteristics. In 2002, a molecular phylogenetic study showed Sundacarpus is embedded in Prumnopitys and the monophyly of Lagarostrobos is doubtful if Manoao is included within it. More recent treatments of the family have recognized Manoao, but not Sundacarpus.
Evolution
Molecular evidence supports Podocarpaceae being the sister group to the Araucariaceae, and having diverged from it during the late Permian. While some fossils attributed to the family have been reported from the Late Permian and Triassic, like Rissikia, these cannot be unambiguously assigned to the family. The oldest unambiguous members of the family are known from the Jurassic period, found across both hemispheres, such as Scarburgia and Harrisiocarpus from the Middle Jurassic of England, as well as unnamed species from the Middle-Late Jurassic of Patagonia. Modern genera of the family first appeared during the Early Cretaceous, with the family probably reaching an apex of diversity during the early Cenozoic.
Genera
Studies based on anatomical, biogeographical, morphological, and DNA evidence suggest these relationships:
List of genera
| Biology and health sciences | Pinophyta (Conifers) | Plants |
457977 | https://en.wikipedia.org/wiki/Micropaleontology | Micropaleontology | Micropaleontology (American spelling; spelled micropalaeontology in European usage) is the branch of paleontology (palaeontology) that studies microfossils, or fossils that require the use of a microscope to see the organism, its morphology and its characteristic details.
Microfossils
Microfossils are fossils that are generally between 0.001mm and 1 mm in size, the study of which requires the use of light or electron microscopy. Fossils which can be studied by the naked eye or low-powered magnification, such as a hand lens, are referred to as macrofossils.
For example, some colonial organisms, such as Bryozoa (especially the Cheilostomata) have relatively large colonies, but are classified by fine skeletal details of the small individuals of the colony.
In another example, many fossil genera of Foraminifera, which are protists are known from shells (called "tests") that were as big as coins, such as the genus Nummulites.
Microfossils are a common feature of the geological record, from the Precambrian to the Holocene. They are most common in deposits of marine environments, but also occur in brackish water, fresh water and terrestrial sedimentary deposits. While every kingdom of life is represented in the microfossil record, the most abundant forms are protist skeletons or cysts from the Chrysophyta, Pyrrhophyta, Sarcodina, acritarchs and chitinozoans, together with pollen and spores from the vascular plants.
In 2017, fossilized microorganisms, or microfossils, were announced to have been discovered in hydrothermal vent precipitates in the Nuvvuagittuq Belt of Quebec, Canada that may be as old as 4.28 billion years old, the oldest record of life on Earth, suggesting "an almost instantaneous emergence of life" (in a geological time-scale sense), after ocean formation 4.41 billion years ago, and not long after the formation of the Earth 4.54 billion years ago. Nonetheless, life may have started even earlier, at nearly 4.5 billion years ago, as claimed by some researchers.
Areas of study
Micropaleontology can be roughly divided into four areas of study on the basis of microfossil composition: (a) calcareous, as in coccoliths and foraminifera, (b) phosphatic, as in the study of some vertebrates, (c) siliceous, as in diatoms and radiolaria, or (d) organic, as in the pollen and spores studied in palynology.
This division reflects differences in the mineralogical and chemical composition of microfossil remains (and therefore in the methods of fossil recovery) rather than any strict taxonomic or ecological distinctions. Most researchers in this field, known as micropaleontologists, are typically specialists in one or more taxonomic groups.
Calcareous microfossils
Calcareous (CaCO3) microfossils include coccoliths, foraminifera, calcareous dinoflagellate cysts, and ostracods (seed shrimp).
Phosphatic microfossils
Phosphatic microfossils include conodonts (tiny oral structures of an extinct chordate group), some scolecodonts ("worm" jaws), shark spines and teeth, and other fish remains (collectively called "ichthyoliths").
Siliceous microfossils
Siliceous microfossils include diatoms, radiolarians, silicoflagellates, ebridians, phytoliths, some scolecodonts ("worm" jaws), and sponge spicules.
Organic microfossils
The study of organic microfossils is called palynology. Organic microfossils include pollen, spores, chitinozoans (thought to be the egg cases of marine invertebrates), scolecodonts ("worm" jaws), acritarchs, dinoflagellate cysts, and fungal remains.
Methods
Sediment or rock samples are collected from either cores or outcrops, and the microfossils they contain are extracted by a variety of physical and chemical laboratory techniques, including sieving, density separation by centrifuge or in heavy liquids, and chemical digestion of the unwanted fraction. The resulting concentrated sample of microfossils is then mounted on a slide for analysis, usually by light microscope. Taxa are then identified and counted. The enormous numbers of microfossils that a small sediment sample can often yield allows the collection of statistically robust datasets which can be subjected to multivariate analysis. A typical microfossil study will involve identification of a few hundred specimens from each sample.
Application of micropaleontology
Microfossils are specially noteworthy for their importance in biostratigraphy. Since microfossils are often extremely abundant, widespread, and quick to appear and disappear from the stratigraphic record, they constitute ideal index fossils from a biostratigraphic perspective. Also, the planktonic and nektonic habits of some microfossils give them the bonus of appearing across a wide range of facies or paleoenvironments, as well as having near-global distribution, making biostratigraphic correlation even more powerful and effective.
Microfossils, particularly from deep-sea sediments, also provide some of the most important records of global environmental change on long, medium or short timescales. Across vast areas of the ocean floor, the shells of planktonic micro-organisms sinking from surface waters provide the dominant source of sediment, and they continuously accumulate (typically at rates of 20–50 million per million years). Study of changes in assemblages of microfossils and changes in their shell chemistry (e.g., oxygen isotope composition) are fundamental to research on climate change in the geological past.
In addition to providing an excellent tool for sedimentary rock dating and for paleoenvironmental reconstruction – heavily used in both petroleum geology and paleoceanography – micropaleontology has also found a number of less orthodox applications, such as its growing role in forensic police investigation or in determining the provenance of archaeological artefacts.
Micropaleontology is also a tool of geoarchaeology used in the archaeological reconstruction of human habitation sites and environments. Changes in the microfossil population abundance in the stratigraphy of current and former water bodies reflect changes in environmental conditions. Naturally occurring ostracods in freshwater bodies are impacted by changes in salinity and pH due to human activities. When correlated with other dating techniques, prehistoric environments can be reconstructed. Work on Lake Tanganyika provided a profile of human-induced environmental changes of a 4,000-year period.
Similar work in the arid American Southwest has provided information on irrigation canals used by prehistoric peoples from 2100 B.C. to 500 B.C. Other archaeological work in arid climates throughout the Americas has incorporated Micropaleontological analysis to build a more complete picture of prehistoric climate and human activity.
| Biology and health sciences | Paleontology | Biology |
457991 | https://en.wikipedia.org/wiki/Paleoecology | Paleoecology | Paleoecology (also spelled palaeoecology) is the study of interactions between organisms and/or interactions between organisms and their environments across geologic timescales. As a discipline, paleoecology interacts with, depends on and informs a variety of fields including paleontology, ecology, climatology and biology.
Paleoecology emerged from the field of paleontology in the 1950s, though paleontologists have conducted paleoecological studies since the creation of paleontology in the 1700s and 1800s. Combining the investigative approach of searching for fossils with the theoretical approach of Charles Darwin and Alexander von Humboldt, paleoecology began as paleontologists began examining both the ancient organisms they discovered and the reconstructed environments in which they lived. Visual depictions of past marine and terrestrial communities have been considered an early form of paleoecology. The term "paleo-ecology" was coined by Frederic Clements in 1916.
Overview of paleoecological approaches
Classic paleoecology uses data from fossils and subfossils to reconstruct the ecosystems of the past. It involves the study of fossil organisms and their associated remains (such as shells, teeth, pollen, and seeds), which can help in the interpretation of their life cycle, living interactions, natural environment, communities, and manner of death and burial. Such interpretations aid the reconstruction of past environments (i.e., paleoenvironments). Paleoecologists have studied the fossil record to try to clarify the relationship animals have to their environment, in part to help understand the current state of biodiversity. They have identified close links between vertebrate taxonomic and ecological diversity, that is, between the diversity of animals and the niches they occupy. Classical paleoecology is a primarily reductionist approach: scientists conduct detailed analysis of relatively small groups of organisms within shorter geologic timeframes.
Evolutionary paleoecology uses data from fossils and other evidence to examine how organisms and their environments change throughout time. Evolutionary paleoecologists take the holistic approach of looking at both organism and environmental change, accounting for physical and chemical changes in the atmosphere, lithosphere and hydrosphere across time. By studying patterns of evolution and extinction in the context of environmental change, evolutionary paleoecologists are able to examine concepts of vulnerability and resilience in species and environments.
Community paleoecology uses statistical analysis to examine the composition and distribution of groups of plants or animals. By quantifying how plants or animals are associated, community paleoecologists are able to investigate the structures of ancient communities of organisms. Advances in technology have helped propel the field, through the use of physical models and computer-based analysis.
Major principles
While the functions and relationships of fossil organisms may not be observed directly (as in ecology), scientists can describe and analyze both individuals and communities over time. To do so, paleoecologists make the following assumptions:
All organisms are adapted and restricted to a particular environment, and are usually adapted to a particular lifestyle.
Essentially all organisms depend on another organism, whether directly or indirectly.
The fossil or physical records are inherently incomplete - the geologic record is selective and some environments are more likely to be preserved than others. Taphonomy, affecting the over- and underrepresentation of fossils, is an extremely important consideration in interpreting fossil assemblages.
Uniformitarianism is the concept that processes that took place in the geologic past are the same as the ones that are observed taking place today. In paleoecology, uniformitarianism is used as a methodology: paleoecologists make inferences about ancient organisms and environments based on analogies they find in the present.
Paleoecological methods
The aim of paleoecology is to build the most detailed model possible of the life environment of previously living organisms found today as fossils. The process of reconstructing past environments requires the use of archives (e.g., sediment sequences), proxies (e.g., the micro or mega-fossils and other sediment characteristics that provide the evidence of the biota and the physical environment), and chronology (e.g., obtaining absolute (or relative) dating of events in the archive). Such reconstruction takes into consideration complex interactions among environmental factors such as temperatures, food supplies, and degree of solar illumination. Often much of this information is lost or distorted by the fossilization process or diagenesis of the enclosing sediments, making interpretation difficult.
Some other proxies for reconstructing past environments include charcoal and pollen, which synthesize fire and vegetation data, respectively. Both of these alternates can be found in lakes and peat settings, and can provide moderate to high resolution information. These are well studied methods often utilized in the paleoecological field.
The environmental complexity factor is normally tackled through statistical analysis of the available numerical data (quantitative paleontology or paleostatistics), while the study of post-mortem processes is known as the field of taphonomy.
Quaternary
Because the Quaternary period is well represented in geographically extensive and high temporal-resolution records, many hypotheses arising from ecological studies of modern environments can be tested at the millennial scale using paleoecological data. In addition, such studies provide historical (pre-industrialization) baselines of species composition and disturbance regimes for ecosystem restoration, or provide examples for understanding the dynamics of ecosystem change through periods of large climate changes. Paleoecological studies are used to inform conservation, management and restoration efforts. In particular, fire-focused paleoecology is an informative field of study to land managers seeking to restore ecosystem fire regimes.
| Biology and health sciences | Paleontology | Biology |
458086 | https://en.wikipedia.org/wiki/Rattlesnake | Rattlesnake | Rattlesnakes are venomous snakes that form the genera Crotalus and Sistrurus of the subfamily Crotalinae (the pit vipers). All rattlesnakes are vipers. Rattlesnakes are predators that live in a wide array of habitats, hunting small animals such as birds and rodents.
Rattlesnakes receive their name from the rattle located at the end of their tails, which makes a loud rattling noise when vibrated that deters predators. Rattlesnakes are the leading contributor to snakebite injuries in North America, but rarely bite unless provoked or threatened; if treated promptly, the bites are seldom fatal.
The 36 known species of rattlesnakes have between 65 and 70 subspecies, all native to the Americas, ranging from central Argentina to southern Canada. The largest rattlesnake, the eastern diamondback, can measure up to in length.
Rattlesnakes are preyed upon by hawks, weasels, kingsnakes, and a variety of other species. Rattlesnakes are heavily preyed upon as neonates, while they are still weak and immature. Large numbers of rattlesnakes are killed by humans. Rattlesnake populations in many areas are severely threatened by habitat destruction, poaching, and extermination campaigns.
Etymology
The scientific name Crotalus is derived from the Greek κρόταλον, meaning "castanet". The name Sistrurus is the Latinized form of the Greek word for "tail rattler" (Σείστρουρος, seistrouros) and shares its root with the ancient Egyptian musical instrument the sistrum, a type of rattle.
Ecology
Range and habitat
Rattlesnakes are native to the Americas from southern Canada to central Argentina, with the majority of species inhabiting arid regions. The large majority of species live in the American Southwest and Mexico. Four species may be found east of the Mississippi River, and two in South America. In the United States, the state with the most types of rattlesnakes is Arizona, with 13.
Rattlesnakes are found in almost every habitat type capable of supporting terrestrial ectothermic vertebrates, but individual species may have extremely specific habitat requirements, living where certain plant associations occur or within a narrow range of elevations. Most species live near open, rocky areas. Rocks offer them cover from predators, plentiful prey (e.g. rodents, lizards, insects, etc. that live amidst the rocks), and open basking areas. However, rattlesnakes can also be found in a wide variety of other habitats, including prairies, marshes, deserts, and forests. Rattlesnakes prefer a temperature range between 80 and 90 °F (26 and 32 °C), but can survive temperatures below freezing, recovering from brief exposure to temperatures as low as 4 °F (−16 °C), and surviving for several days in temperatures as low as 37 °F (3 °C).
The most probable ancestral area of rattlesnakes is the Sierra Madre Occidental region in Mexico. The most probable vegetation or habitat of the ancestral area appears to be pine-oak forests.
Prey
Rattlesnakes typically consume mice, rats, rabbits, squirrels, small birds, and other small animals. They lie in wait for their prey, or hunt for it in holes. The rattlesnake's defence and hunting mechanisms are bound to its physiology and its environment. More importantly environmental temperature can influence the ability of ectotherms. The prey is killed quickly with a venomous bite as opposed to constriction. If the bitten prey moves away before dying, the rattlesnake can follow it by its scent. When it locates the fallen prey, it checks for signs of life by prodding with its snout, flicking its tongue, and using its sense of smell. Once the prey has become incapacitated, the rattlesnake locates its head by odors emitted from the mouth. The prey is then ingested head first, which allows wings and limbs to fold at the joints in a manner that minimizes the girth of the meal. The gastric fluids of rattlesnakes are extremely powerful, allowing for the digestion of flesh and bone. Optimal digestion occurs when the snake maintains a body temperature between 80 and 85 °F (25 and 29 °C). If the prey is small, the rattlesnake often continues hunting. If the meal was adequate, the snake finds a warm, safe location in which to coil up and rest until the prey is digested. Feeding habits play an important ecological role by limiting the size of rodent populations, which prevents crop damage and stabilizes ecosystems.
Hydration
Rattlesnakes are believed to require at least their own body weight in water annually to remain hydrated. The method by which they drink depends on the water source. In larger bodies of water (streams, ponds, etc.), they submerge their heads and ingest water by opening and closing their jaws, which sucks in water. If drinking dew or small puddles, they sip the liquid either by capillary action or by flattening and flooding their lower jaws.
In desert environments, scientists observe that rattlesnakes have evolved to stay hydrated by coiling up and flattening so that their bodies can collect rain. In some instances, they aggregated in a carpet-like formation to make a larger rain-collecting platform.
Predators
Newborn rattlesnakes are heavily preyed upon by a variety of species, including cats, ravens, crows, roadrunners, raccoons, opossums, skunks, coyotes, weasels, whipsnakes, kingsnakes, and racers. Young of the smaller crotaline species are frequently killed and eaten by small predatory birds, such as jays, kingfishers, and shrikes. Some species of ants in the genus Formica are known to prey upon neonates, and Solenopsis invicta (fire ants) likely do, as well. On occasion, hungry adult rattlesnakes cannibalize neonates. The small proportion (often as few as 20%) of rattlesnakes that make it to their second year are heavily preyed upon by a variety of larger predators, including coyotes, eagles, hawks, owls, falcons, feral pigs, badgers, indigo snakes, and kingsnakes.
The common kingsnake (Lampropeltis getula), a constrictor, is immune to the venom of rattlesnakes and other vipers, and rattlesnakes form part of its natural diet. Rattlesnakes sense kingsnakes' presence by their odor. When they realize a kingsnake is nearby, they begin enacting a set of defensive postures known as "body bridging". Unlike its normal erect and coiled defensive-striking posture, the rattlesnake keeps its head low to the ground in an attempt to prevent the kingsnake from gaining a hold on it (the head being the first part of the rattlesnake to be ingested). The rattlesnake jerks its body about, while bridging its back upwards, forming an elevated coil that faces the kingsnake. The elevated coil is used to strike the attacker, and is also used to shield the head from the kingsnake.
Anatomy
Sensory organs
Like all pit vipers, rattlesnakes have two organs that can sense radiation; their eyes and a set of heat-sensing "pits" on their faces that enable them to locate prey and move towards it, based on the prey's thermal radiation signature. These pits have a relatively short effective range of about but give the rattlesnake a distinct advantage in hunting for warm-blooded creatures at night.
Heat-sensing pits
Aside from their eyes, rattlesnakes are able to detect thermal radiation emitted by warm-blooded organisms in their environment. Functioning optically like a pinhole camera eye, thermal radiation in the form of infrared light passes through the opening of the pit and strikes the pit membrane located in the back wall, warming this part of the organ. Due to the high density of heat-sensitive receptors innervating this membrane, the rattlesnake can detect temperature changes of 0.003 °C or less in its immediate surroundings. Infrared cues from these receptors are transmitted to the brain by the trigeminal nerve, where they are used to create thermal maps of the snake's surroundings. Due to the small sizes of the pit openings, typically these thermal images are low in resolution and contrast. Nevertheless, rattlesnakes superimpose visual images created from information from the eyes with these thermal images from the pit organs to more accurately visualize their surroundings in low levels of light. Research conducted recently on the molecular mechanism of this ability suggests the temperature sensitivity of these pit organs is closely linked to the activity of transient receptor potential ankyrin 1, a temperature-sensitive ion channel saturated in the pit membrane.
Eyes
Rattlesnake eyes, which contain many rod cells, are well adapted to nocturnal use. Rattlesnakes, though, are not exclusively nocturnal, and their vision is more acute during daylight conditions. Rattlesnakes also possess cone cells, which means they are capable of some form of color vision. The rattlesnake eye lacks a fovea, making vision of sharply defined images impossible. Instead, they mostly rely on the perception of movement. Rattlesnake eyes are capable of horizontal rotation, but they do not appear to move their eyeballs to follow moving objects.
Smell
Rattlesnakes have an exceptionally keen sense of smell. They can sense olfactory stimuli both through their nostrils and by flicking their tongues, which carry scent-bearing particles to the Jacobson's organs in the roof of their mouths.
Auditory system
Like all snakes, rattlesnakes lack external ear openings, and the structures of their middle ear are not as highly specialized as those of other vertebrates, such as mammals. Thus, their sense of hearing is not very effective, but they are capable of sensing vibrations in the ground, passed by the skeleton to the auditory nerve.
Fangs
Rattlesnake fangs are connected by venom ducts to large venom glands near the outer edge of the upper jaw, towards the rear of the head. When the rattlesnake bites, muscles on the sides of the venom glands contract to squeeze the venom through the ducts and into the fangs. When the fangs are not in use, they remain folded against the palate.
Rattlesnakes are born with fully functioning fangs and venom, and are capable of killing prey at birth. Adult rattlesnakes shed their fangs every 6–10 weeks. At least three pairs of replacement fangs lie behind the functional pair.
Venom
Rattlesnake venom is hemotoxic, destroying tissue, causing necrosis and coagulopathy (disrupted blood clotting). In the U.S., the tiger rattlesnake (C. tigris) and some varieties of the Mojave rattlesnake (C. scutulatus) also have a presynaptic neurotoxic venom component known as Mojave type A toxin, which can cause severe paralysis. However, most North American rattlesnakes are not neurotoxic. Although it has a comparatively low venom yield, the venom toxicity of C. tigris is considered to be among the highest of all rattlesnake venoms, and among the highest of all snakes in the Western Hemisphere based on studies conducted on laboratory mice. C. scutulatus is also widely regarded as producing one of the most toxic snake venoms in the Americas, based on studies in laboratory mice.
Rattlesnake venom is a mixture of five to fifteen enzymes, various metal ions, biogenic amines, lipids, free amino acids, proteins, and polypeptides. More specifically, there are three main families of toxins in rattlesnakes: phospholipases A2 (PLA2s), snake venom metalloproteinases (SVMPs), and snake venom serine proteinases (SVSPs). It contains components meant to immobilize and disable the prey, as well as digestive enzymes, which break down tissue to prepare for later ingestion. The venom is very stable, and retains its toxicity for many years in storage.
Snake venom, in general, has a complex and ongoing evolutionary process, and rattlesnake venom is no different. The primary mechanisms of evolution are both gene duplication and gene loss events. The duplication events provided material for neofunctionalization to create the novel toxin genes, while gene loss influenced speciation and helped lead to such a wide variety of “chemical cocktails” in rattlesnake venoms. The prevailing theory for the driving force of this evolution is directional selection, where efficacy on prey is selected for. Diversity in prey leads to less specificity in toxins, while highly specialized toxins are more likely to develop when there are few key prey species. However, recently, balancing selection has been indicated to better explain the maintenance of adaptive genetic diversity in venom-related genes, potentially allowing for the rattlesnakes to better keep up in the evolutionary arms race with their prey.
Older snakes possess more potent venom, and larger snakes are frequently capable of storing larger volumes of it.
Rattle
The rattle serves as a warning for predators of the rattlesnake. The rattle is composed of a series of hollow, interlocked segments made of keratin, which are created by modifying the scales that cover the tip of the tail. The contraction of special "shaker" muscles in the tail causes these segments to vibrate against one another, thus making the rattling noise (which is amplified because the segments are hollow) in a behavior known as tail vibration. The muscles which cause rattling are some of the fastest known, firing 50 times per second on average, sustainable for a duration of up to three hours.
In 2016, Allf et al. published a paper proposing behavioral plasticity as the mechanism by which the rattling system evolved in rattlesnakes. In the case of rattlesnakes, Allf et al. proposed that tail vibration in response to predator threat could be the precursor for the rattling system in rattlesnakes, an example of behavioral plasticity. To investigate this hypothesis, the researchers analyzed tail vibration and relatedness to rattlesnakes among snakes in the families Viperidae and Colubridae. Their results demonstrated that the more closely related a species was to rattlesnakes, the more similar that species was to rattlesnakes in both duration and rate of tail vibration. These results strongly support the hypothesis that tail vibration preceded the rattling system as a behavior and allowed for the rattle to be selected for once developed. Even a small, underdeveloped rattle early in the evolution of the rattling system could have been advantageous if tail vibration was an ancestral behavior.
At birth, a "prebutton" is present at the tip of the snake's tail; it is replaced by the "button" several days later when the first skin is shed. However, no sound can be made by the rattle until a second segment is added when the skin is shed again. A new rattle segment is added each time the snake sheds its skin, and the snake may shed its skin several times a year, depending on food supply and growth rate.
Rattlesnakes travel with their rattles held up to protect them from damage, but in spite of this precaution, their day-to-day activities in the wild still cause them to regularly break off end segments. Because of this, the number of rattles on its tail is not related to the age of a rattlesnake.
Compared to females, males have thicker and longer tails (because they contain the inverted hemipenes). Also, the tails of males taper gradually from the body, whereas the tails of females narrow abruptly at the vent.
Skin and circulation
Rattlesnakes, like other members of the Squamata order, have a circulatory system that is powered by a three-chambered heart composed of two atria and one ventricle. The right atrium receives deoxygenated blood from veins coming from the systemic circuit. The left atrium receives oxygenated blood from the lungs in the pulmonary circuit and pumps it to the ventricle and through the systemic circuit via capillaries and arteries.
Rattlesnake skin has a set of overlapping scales that cover the entire body, providing protection from a variety of threats, including dehydration and physical trauma. The typical rattlesnake, genus Crotalus, has the top of its head covered with small scales, except, with a few species, a few crowded plates directly over the snout. The skin of snakes is highly sensitive to contact, tension, and pressure; they are capable of feeling pain.
An important function of the skin is the sensation of changes in air temperature, which can guide the snakes towards warm basking/shelter locations. All snakes are ectotherms. To maintain a stable body temperature, they exchange heat with their external environments. Snakes often move into open, sunny areas to absorb heat from the sun and warmed earth, a behavior known as basking. Nerves in the skin regulate the flow of blood into the veins near the surface.
The skin of rattlesnakes is intricately patterned in a manner that camouflages them from their predators. Rattlesnakes do not generally have bright or showy colors (reds, yellows, blues, etc.), instead relying on subtle earth tones that resemble the surrounding environment.
Creases in the epidermal tissue connect the scales of rattlesnakes. When ingesting large prey, these creases can unfold, allowing the skin to expand to envelop a much greater volume. The skin appears to tightly stretch to accommodate the meal, but in reality, the skin is simply smoothing out from its creased state and is not under very high tension.
Reproduction
Most rattlesnake species mate during the summer or fall, while some species mate only in the spring, or during both the spring and fall.
Females secrete small amounts of sex pheromones, which leave a trail the males follow using their tongues and Jacobson's organs as guides. Once a receptive female has been located, the male often spends several days following her around (a behavior not common outside of the mating season), frequently touching and rubbing her in an attempt to stimulate her.
The males of some species, such as timber rattlesnakes (C. horridus), fight each other during the mating season, in competition over females. These fights, known as "combat dances", consist of the two males intertwining the anterior portion of their bodies, often with their heads and necks held vertically. The larger males usually end up driving the smaller males away.
Although many kinds of snakes and other reptiles are oviparous (lay eggs), rattlesnakes are ovoviviparous (give birth to live young after carrying eggs inside). The female produces the ova ("eggs") in her ovaries, after which they pass through her body cavity and into one of her two oviducts. The ova are arranged in a continuous chain in a coiled section of the oviduct, known as the "tuba". Male rattlesnakes have sexual organs known as hemipenes, located in the base of the tail. The hemipenis is retracted inside of the body when mating is not occurring. The hemipenis is similar to the human penis. Females can store semen for months in internal recesses known as spermathecae, which permits them to mate during the fall, but not fertilize the ova until the following spring. The Arizona black rattlesnake (C. oreganus cerberus), has been observed to exhibit complex social behavior reminiscent of that in mammals. Females often remain with their young in nests for several weeks, and mothers have been observed co-operatively parenting their broods.
Rattlesnakes generally take several years to mature, and females usually reproduce only once every three years.
Brumation
In the colder winter months, some rattlesnake species enter a period of brumation, which is dormancy similar to hibernation.
They often gather together for brumation in large numbers (sometimes over 1,000 snakes), huddling together inside underground "rattlesnake dens" or hibernacula. They regularly share their winter burrows with a wide variety of other species (such as turtles, small mammals, invertebrates, and other types of snakes).
Rattlesnakes often return to the same den, year after year, sometimes traveling several miles to get there. How the rattlesnakes find their way back to the dens each year is unknown, but may involve a combination of pheromone trails and visual cues (e.g., topography, celestial navigation, and solar orientation).
Species with long periods of brumation tend to have much lower reproductive rates than those with shorter brumation periods or those that do not brumate at all. Female timber rattlesnakes in high peaks in the Appalachian Mountains of New England reproduce every three years on average; the lance-headed rattlesnake (C. polystictus), native to the warm climate of Mexico, reproduces annually.
Like most other snakes, rattlesnakes aestivate during very hot or dry periods, which is why they are rarely seen during the hottest and driest months of summer.
Conservation status
Rattlesnakes tend to avoid developed areas, preferring undisturbed, natural habitats. Rapid habitat destruction by humans, mass killings during events such as rattlesnake round-ups, and deliberate extermination campaigns all pose threats to rattlesnake populations in many areas. Several species, such as the timber rattlesnake, massasauga, and canebrake rattlesnake, are listed as threatened or endangered in many U.S. states.
Many rattlesnakes die from being run over by cars.
In more heavily populated and trafficked areas, reports have been increasing of rattlesnakes that do not rattle. This phenomenon is falsely attributed to selective pressure by humans, who often kill the snakes when they are discovered. However, snake experts have dismissed this theory, stating that snakes simply do not rattle as often as laymen expect them to, and that snakes that live near populated areas simply get used to people passing by, only rattling when a person lingers too long or gets too close.
Safety and first aid
Rattlesnakes are the leading cause of snakebite injuries in North America and a significant cause in Central and South America.
Avoiding bites
Rattlesnakes tend to avoid wide-open spaces where they cannot hide from predators, and generally avoid humans if they are aware of their approach. Rattlesnakes rarely bite unless they feel threatened or provoked. A majority of victims (about 72%) are males. Around half of bites occur in cases where the victim saw the snake, yet made no effort to move away.
Harassing or attacking a rattlesnake, illegal in some jurisdictions, puts one at much higher risk of a bite. Rattlesnakes seek to avoid humans and other predators or large herbivores that themselves pose lethal danger. Dogs, often much more aggressive than humans, are much more likely to experience a snakebite, and are more likely to die of a rattlesnake bite although they can be vaccinated against them.
Caution is advised even when snakes are believed to be dead; rattlesnake heads can sense, flick the tongue, and inflict venomous bites reflexively for up to an hour after being severed from the body.
Effect of bites on humans
An estimated 7,000 to 8,000 people are bitten by venomous snakes in the United States each year, with about five deaths. The most important factor in survival following a severe envenomation is the time elapsed between the bite and treatment. Most deaths occur between 6 and 48 hours after the bite. If antivenom treatment is given within two hours of the bite, the probability of recovery is greater than 99%.
When a bite occurs, the amount of venom injected is under voluntary control by the snake. The amount released depends on a variety of factors, including the condition of the snake (e.g., having long, healthy fangs and a full venom sack) and its temperament (an angry, hungry snake that has just been stepped on vs. a satiated snake that was merely surprised by walking near it). About 20% of bites result in no envenomation at all. A lack of burning pain and edema away from the fang marks after one hour suggests either no or minimal envenomation occurred. A lack of edema or erythema in the area of the bite after eight hours indicates a lack of envenomation for most rattlesnake bites.
Common symptoms include swelling, severe pain, tingling, weakness, anxiety, nausea and vomiting, hemorrhaging, perspiration, and (rarely) heart failure. Local pain following envenomation is often intense, increasing with the ensuing edema. Children generally experience more severe symptoms because they receive a larger amount of venom per unit of body mass.
Antivenom
Antivenom, or antivenin, is commonly used to treat the effects of local and systemic pit viper envenomations. The first step in the production of crotaline antivenom is collecting ("milking") the venom of a live rattlesnake—usually from the western diamondback (Crotalus atrox), eastern diamondback (Crotalus adamanteus), South American rattlesnake (Crotalus durissis terrificus), or fer-de-lance (Bothrops atrox). The extracted venom is then diluted and injected into horses, goats, or sheep, whose immune systems produce antibodies that protect from the toxic effects of the venom. These antibodies accumulate in the blood, which is then extracted and centrifuged to separate the red blood cells. The resulting serum is purified into a lyophilized powder, which is packaged for distribution and later use by human patients.
Because antivenom is derived from animal antibodies, people generally display an allergic response during infusion, known as serum sickness.
Veterinary care
In the United States, more than 15,000 domesticated animals are bitten by snakes each year. Rattlesnake envenomations account for 80% of the deadly incidents.
Dogs are most commonly bitten on the front legs and head. Horses generally receive bites on the muzzle, and cattle on their tongues and muzzles. If a domesticated animal is bitten, the hair around the bite should be removed so the wound can be clearly seen. The crotaline Fab antivenom has been shown to be effective in the treatment of canine rattlesnake bites. Symptoms include swelling, slight bleeding, sensitivity, shaking, and anxiety.
In human culture
Spirituality
Indigenous Americans
Aztec paintings, Central American temples, and the great burial mounds in the Southern United States are frequently adorned with depictions of rattlesnakes, often within the symbols and emblems of the most powerful deities.
The Feathered Serpent of Mesoamerican religion was depicted as having the combined features of the quetzal and rattlesnake. The Ancient Maya considered the rattlesnake to be a "vision serpent" that acted as a conduit to the "otherworld".
Rattlesnakes are a key element in Aztec mythology and were widely represented in Aztec art, including sculptures, jewelry, and architectural elements.
Christian snake-handling sects
Members of some Christian sects in the Southern United States are regularly bitten while participating in "snake handling" rituals. Snake handling is when people hold venomous snakes, unprotected, as part of a religious service inspired by a literal interpretation of the Bible verses Mark 16:17–18, which reads, "In my name ... They shall take up serpents ...."
In traditional medicine
Meriwether Lewis of the Lewis and Clark Expedition described in 1805 how a hired interpreter, who lived for 15 years with the Mandan, used the rattlesnake‘s rattle to speed up the delivery of Sacagewea‘s son Jean Baptiste Charbonneau:
As food
Journalist Alistair Cooke claimed that rattlesnake tastes "just like chicken, only tougher". Others have compared the flavor to a wide range of other meats, including veal, frog, tortoise, quail, fish, rabbit, and even canned tuna. Methods of preparation include barbecueing and frying; author Maud Newton, following a recipe by Harry Crews, described the taste, "at least when breaded and fried, like a sinewy, half-starved tilapia."
Symbolism
The rattlesnake became a symbolic animal for the Colonials during the Revolutionary War period, and is depicted prominently on the Gadsden Flag. It continues to be used as a symbol by the United States military, and political movements within the United States.
| Biology and health sciences | Reptiles | null |
458150 | https://en.wikipedia.org/wiki/Monoplacophora | Monoplacophora | Monoplacophora , meaning "bearing one plate", is a polyphyletic superclass of molluscs with a cap-like shell, inhabiting deep sea environments. Extant representatives were not recognized as such until 1952; previously they were known only from the fossil record, and were thought to have become extinct 375 million years ago.
Although the shell of many monoplacophorans is limpet-like in shape, they are not gastropods, nor do they have any close relation to gastropods.
Definition
Discussion about monoplacophorans is made difficult by the slippery definition of the taxon; some authors take it to refer to all non-gastropod molluscs with a single shell, or all single-shelled molluscs with serially repeated units; whereas other workers restrict the definition to cap-shaped forms, excluding spiral and other shapes of shell. The inclusion of the gastropod-like Bellerophontoidea within the group is also contentious.
One attempt to resolve this confusion was to separate out the predominantly coiled helcionelloids from the traditional, cap-like tergomyans, this latter group containing extant Tryblidiids.
Taxonomy
Taxonomy of Monoplacophora per Bouchet, et al. (2017):
Class Monoplacophora
subclass Cephalopoda?
† Subclass Cyrtolitiones
† Order Sinuitopsida
† Superfamily Cyrtolitoidea S. A. Miller, 1889
† Family Cyrtolitidae S. A. Miller, 1889
† Family Carcassonnellidae Horný, 1997
† Superfamily Cyclocyrtonelloidea Horný, 1962
† Family Cyclocyrtonellidae Horný, 1962 (= Yochelsoniidae Horný, 1962 (inv.))
† Family Multifariitidae Bjaly, 1973
† Family Sinuellidae Starobogatov & Moskalev, 1987
† Family Sinuitinidae Starobogatov & Moskalev, 1987
† Subclass Cyrtonelliones
† Order Cyrtonellida
† Superfamily Cyrtonelloidea Knight & Yochelson, 1958
† Family Cyrtonellidae Knight & Yochelson, 1958 (= Cyrtonellopsinae Horný, 1965)
† Subclass Eomonoplacophora
(Unassigned to Order)
† Superfamily Maikhanelloidea Missarzhevsky, 1989
† Family Maikhanellidae Missarzhevsky, 1989 (= Purellidae Vassiljeva, 1990)
Subclass Tergomya (= Pilinea)
† Order Kirengellida (= Romaniellida)
† Superfamily Archaeophialoidea Knight & Yochelson, 1958
† Family Archaeophialidae Knight & Yochelson, 1958
† Family Peelipilinidae Horný, 2006
† Family Pygmaeoconidae Horný, 2006
† Superfamily Kirengelloidea Starobogatov, 1970
† Family Kirengellidae Starobogatov, 1970
† Family Romaniellidae Rozov, 1975
† Family Nyuellidae Starobogatov & Moskalev, 1987
† Superfamily Hypseloconoidea Knight, 1952
† Family Hypseloconidae Knight, 1952
Order Tryblidiida
Superfamily Tryblidioidea Pilsbry, 1899
† Family Tryblidiidae Pilsbry, 1899
† Family Proplinidae Knight & Yochelson, 1958
† Family Drahomiridae Knight & Yochelson, 1958
† Family Bipulvinidae Starobogatov, 1970
Superfamily Neopilinoidea Knight & Yochelson, 1958
Family Neopilinidae Knight & Yochelson, 1958
Subfamily Neopilininae Knight & Yochelson, 1958 (= Vemidae Moskalev, Starobogatov & Filatova, 1983; = Laevipilinidae Moskalev, Starobogatov & Filatova, 1983; = Monoplacophoridae Moskalev, Starobogatov & Filatova, 1983)
Subfamily Veleropilininae Starobogatov & Moskalev, 1987 (= Rokopellidae Starobogatov & Moskalev, 1987; = Micropilinidae Haszprunar & Schaefer, 1997)
Anatomy and physiology
Monoplacophorans are univalved (though not gastropodal), limpet-shaped, and are untorted. They have a pseudometamerism of bilaterally symmetrical repeated organs and muscles. The extant members of the class live only in the deep ocean (the abyssal zone, the continental shelf, and the continental slope) at depths below . Cambrian forms predominately lived in shallow seas, whereas later Paleozoic forms are more commonly found in deeper waters with soft, muddy sea floors.
Although superficially resembling limpets when viewed dorsally, monoplacophorans are not anatomically similar to gastropods. Some similarities are shared with the chitons, such as having segmented anatomy (organs arranged in series). There are eight pairs of dorso-ventral muscles (shell muscles). The nervous system is relatively simple, with no true ganglion present.
The repeated organs include from three to six pairs of "gills" (actually ctenidia) located in a curved line along each side of the foot (though the number is not always considered definitive of a given species), and as many as six "kidneys" (actually nephridia). The tip or point of their low shells points forward rather than towards the back. The shell ranges from 3 mm to 37 mm in diameter depending on species. Like in chitons, the head is poorly defined, and there are no eyes. The mouth is located within the animal's undeveloped head in front of its single large foot and contains a radula, a defining characteristic of the mollusca. Tentacles are situated behind the mouth. They also have a cone-shaped stomach with a single crystalline style though no gastric shield. The intestines are long and make between four and six loops before reaching the posteriorly-positioned anus. Monoplacophorans also have oesophageal pouches.
The sexes are separate with any given animal having two pairs of either ovaries or testes connected to either the third or fourth pair of kidneys. One genus, Micropilina, has apparently been recorded as brooding young in the distal oviduct and pallial groove, releasing the young when approximately 300 micrometers in diameter.
Phylogenetic position
In 2006 a molecular study on Laevipilina antarctica suggested that extant Monoplacophora and Polyplacophora form a well-supported clade with the researched Neopilina closest to the chitons. The two classes in this new clade, with the proposed name Serialia, all show a variable number of serially repeated gills and eight sets of dorsoventral pedal retractor muscles.
This study contradicts the fossil evidence, which suggests that the Monoplacophora are the sister group to the remainder of the conchiferans, and that the cephalopods (squids, octopuses, and relatives) arose from within the monoplacophoran lineage. However, some authors dispute this view and do not necessarily see modern Monoplacophora as related to their presumed fossil ancestors.
The concept of Serialia is supported by other molecular studies.
The fossil record does indicate that the ancestral mollusc was monoplacophoran-like and that the Polyplacophora arose from within the Monoplacophora – not the other way around. This could be reconciled if a secondary loss of shells caused a monoplacophoran body form to re-appear secondarily, which is plausible: At the very least, modern monoplacophorans are not closely related to vent-dwelling representatives from the Silurian.
Cambrian monoplacophoran Knightoconus antarcticus is hypothesised to be an ancestor to the cephalopods.
Fossil species
Living families:
Tryblidiida
Laevipilinidae
Micropilinidae
Monoplacophoridae
Neopilinidae
Extinct families:
Tryblidiida
† Tryblidiidae von Zittel, 1899
† Palaeacmaeidae (uncertain, as the Cambrian type species is a cnidarian. It is maintained here as a receptacle for the Paleozoic genus Parmophorella)
† Palaeacmaea Hall & Whitfield, 1872
† Parmophorella Matthew, 1886
† Knightoconus
† Knightoconus antarcticus
Many Cambrian-Devonian species have been described as "monoplacophorans", but the only fossil members of the crown group date to the Pleistocene.
The Taxonomy of the Gastropoda (Bouchet & Rocroi, 2005) also contains Paleozoic molluscs of uncertain systematic position. It is not known whether these were gastropods or monoplacophorans.
| Biology and health sciences | Mollusks | Animals |
458202 | https://en.wikipedia.org/wiki/Semiheavy%20water | Semiheavy water | Semiheavy water is the result of replacing one of the protium (normal hydrogen, H) in normal water with deuterium (H; or less correctly, D). It exists whenever there is water with H and H in the mix. This is because hydrogen atoms (H) are rapidly exchanged between water molecules. Water with 50% H and 50% H, is about 50% HHO and 25% each of HO and HO, in dynamic equilibrium.
In normal water, about 1 molecule in 3,200 is HDO (HHO) (one hydrogen in 6,400 is H). By comparison, heavy water DO or HO occurs at a proportion of about 1 molecule in 41 million (i.e., 1 in 6,400). This makes semiheavy water far more common than "normal" heavy water.
The freezing point of semiheavy water is close to the freezing point of heavy water at 3.8°C compared to the 3.82°C of heavy water.
Production
On Earth, semiheavy water occurs naturally in normal water at a proportion of about 1 molecule in 3,200; because 1 in 6,400 hydrogen atoms in water is deuterium, which is 1 part in 3,200 by weight. HDO may be separated from normal water by distillation or electrolysis, or by various chemical exchange processes, all of which exploit a kinetic isotope effect. Partial enrichment also occurs in natural bodies of water under certain evaporation conditions. (For more information about the distribution of deuterium in water, see Vienna Standard Mean Ocean Water and Hydrogen isotope biogeochemistry.)
| Physical sciences | Water | Chemistry |
458350 | https://en.wikipedia.org/wiki/Azurite | Azurite | Azurite or Azure spar is a soft, deep-blue copper mineral produced by weathering of copper ore deposits. During the early 19th century, it was also known as chessylite, after the type locality at Chessy-les-Mines near Lyon, France. The mineral, a basic carbonate with the chemical formula Cu3(CO3)2(OH)2, has been known since ancient times, and was mentioned in Pliny the Elder's Natural History under the Greek name (κυανός: "deep blue," root of English cyan) and the Latin name caeruleum. Copper (Cu2+) gives it its blue color.
Mineralogy
Azurite has the formula Cu3(CO3)2(OH)2, with the copper(II) cations linked to two different anions, carbonate and hydroxide. It is one of two relatively common basic copper(II) carbonate minerals, the other being bright green malachite. Aurichalcite is a rare basic carbonate of copper and zinc. Simple copper carbonate (CuCO3) is not known to exist in nature, due to the high affinity of the ion for the hydroxide anion .
Azurite crystallizes in the monoclinic system. Large crystals are dark blue, often prismatic. Azurite specimens can be massive to nodular or can occur as drusy crystals lining a cavity.
Azurite has a Mohs hardness of 3.5 to 4. The specific gravity of azurite is 3.7 to 3.9. Characteristic of a carbonate, specimens effervesce upon treatment with hydrochloric acid. The combination of deep blue color and effervescence when moistened with hydrochloric acid are identifying characteristics of the mineral.
Color
The optical properties (color, intensity) of minerals such as azurite and malachite are characteristic of copper(II). Many coordination complexes of copper(II) exhibit similar colors. According to crystal field theory, the color results from low energy d-d transitions associated with the d9 metal center.
Weathering
Azurite is unstable in open air compared to malachite, and often is pseudomorphically replaced by malachite. This weathering process involves the replacement of some of the carbon dioxide (CO2) units with water (H2O), changing the carbonate:hydroxide ratio of azurite from 1:1 to the 1:2 ratio of malachite:
2 Cu3(CO3)2(OH)2 + H2O → 3 Cu2(CO3)(OH)2 + CO2
From the above equation, the conversion of azurite into malachite is attributable to the low partial pressure of carbon dioxide in air.
Azurite is quite stable under ordinary storage conditions, so that specimens retain their deep blue color for long periods of time.
Occurrences
Azurite is found in the same geologic settings as its sister mineral, malachite, though it is usually less abundant. Both minerals occur widely as supergene copper minerals, formed in the oxidized zone of copper ore deposits. Here they are associated with cuprite, native copper, and various iron oxide minerals.
Fine specimens can be found at many locations. Among the best specimens are found at Bisbee, Arizona, and nearby locations, and have included clusters of crystals several inches long and spherical aggregates and rosettes up to in diameter. Similar rosettes are found at Chessy, Rhône, France. The best crystals, up to in length, are found at Tsumeb, Namibia. Other notable occurrences are in Utah; Mexico; the Ural and Altai Mountains; Sardinia; Laurion, Greece; Wallaroo, South Australia; Brazil and Broken Hill.
Uses
Pigments
Azurite is unstable in air, however it was used as a blue pigment in antiquity. Azurite is naturally occurring in Sinai and the Eastern Desert of Egypt. It was reported by F. C. J. Spurrell (1895) in the following examples; a shell used as a pallet in a Fourth Dynasty (2613 to 2494 BCE) context in Meidum, a cloth over the face of a Fifth Dynasty (2494 to 2345 BCE) mummy also at Meidum and a number of Eighteenth Dynasty (1543–1292 BCE) wall paintings. Depending on the degree of fineness to which it was ground, and its basic content of copper carbonate, it gave a wide range of blues. It has been known as mountain blue, Armenian stone, and azurro della Magna (blue from Germany in Italian). When mixed with oil it turns slightly green. When mixed with egg yolk it turns green-grey. It is also known by the names blue bice and blue verditer, though verditer usually refers to a pigment made by chemical process. Older examples of azurite pigment may show a more greenish tint due to weathering into malachite. Much azurite was mislabeled lapis lazuli, a term applied to many blue pigments. As chemical analysis of paintings from the Middle Ages improves, azurite is being recognized as a major source of the blues used by medieval painters. Lapis lazuli (the pigment ultramarine) was chiefly supplied from Afghanistan during the Middle Ages, whereas azurite was a common mineral in Europe at the time. Sizable deposits were found near Lyons, France. It was mined since the 12th century in Saxony, in the silver mines located there.
Heating can be used to distinguish azurite from purified natural ultramarine blue, a more expensive but more stable blue pigment, as described by Cennino D'Andrea Cennini. Ultramarine withstands heat, whereas azurite converts to black copper oxide. However, gentle heating of azurite produces a deep blue pigment used in Japanese painting techniques.
Azurite pigment can be synthesized by precipitating copper(II) hydroxide from a solution of copper(II) chloride with lime (calcium hydroxide) and treating the precipitate with a concentrated solution of potassium carbonate and lime. This pigment is likely to contain traces of basic copper(II) chlorides.
Jewelry
Azurite is used occasionally as beads and as jewelry, and also as an ornamental stone. However, its softness and tendency to lose its deep blue color as it weathers leaves it with fewer uses. Heating destroys azurite easily, so all mounting of azurite specimens must be done at room temperature.
Collecting
The intense color of azurite makes it a popular collector's stone. The notion that specimens must be carefully protected from bright light, heat, and open air to retain their intensity of color over time may be an urban legend. Paul E. Desautels, former curator of gems and minerals at the Smithsonian Institution, has written that azurite is stable under ordinary storage conditions.
Prospecting
While not a major ore of copper itself, the presence of azurite is a good surface indicator of the presence of weathered copper sulfide ores. It is usually found in association with the chemically similar malachite, producing a striking color combination of deep blue and bright green that is strongly indicative of the presence of copper ores.
History
Azurite was known in the pre-classical ancient world. It was used in ancient Egypt as a pigment, obtained from mines in Sinai. Ancient Mesopotamian writers report the use of a special mortar and pestle for grinding it. It was also used in ancient Greece, for example on the Acropolis in Athens. It does not appear to have been used in ancient Roman wall paintings but Roman writers certainly knew about its use as a pigment. The fusing of glass and azurite was developed in ancient Mesopotamia.
Gallery
| Physical sciences | Minerals | Earth science |
458673 | https://en.wikipedia.org/wiki/Chlorine%20dioxide | Chlorine dioxide | Chlorine dioxide is a chemical compound with the formula ClO2 that exists as yellowish-green gas above 11 °C, a reddish-brown liquid between 11 °C and −59 °C, and as bright orange crystals below −59 °C. It is usually handled as an aqueous solution. It is commonly used as a bleach. More recent developments have extended its applications in food processing and as a disinfectant.
Structure and bonding
The molecule ClO2 has an odd number of valence electrons, and therefore, it is a paramagnetic radical. It is an unusual "example of an odd-electron molecule stable toward dimerization" (nitric oxide being another example).
ClO2 crystallizes in the orthorhombic Pbca space group.
History
In 1933, Lawrence O. Brockway, a graduate student of Linus Pauling, proposed a structure that involved a three-electron bond and two single bonds. However, Pauling in his General Chemistry shows a double bond to one oxygen and a single bond plus a three-electron bond to the other. The valence bond structure would be represented as the resonance hybrid depicted by Pauling. The three-electron bond represents a bond that is weaker than the double bond. In molecular orbital theory this idea is commonplace if the third electron is placed in an anti-bonding orbital. Later work has confirmed that the highest occupied molecular orbital is indeed an incompletely-filled antibonding orbital.
Preparation
Chlorine dioxide was first prepared in 1811 by Sir Humphry Davy.
The reaction of chlorine with oxygen under conditions of flash photolysis in the presence of ultraviolet light results in trace amounts of chlorine dioxide formation.
Cl2 + 2 O2 ->[\ce{UV}] 2 ClO2 ^.
Chlorine dioxide can decompose violently when separated from diluting substances. As a result, preparation methods that involve producing solutions of it without going through a gas-phase stage are often preferred.
Oxidation of chlorite
In the laboratory, ClO2 can be prepared by oxidation of sodium chlorite with chlorine:
Traditionally, chlorine dioxide for disinfection applications has been made from sodium chlorite or the sodium chlorite–hypochlorite method:
or the sodium chlorite–hydrochloric acid method:
or the chlorite–sulfuric acid method:
All three methods can produce chlorine dioxide with high chlorite conversion yield. Unlike the other processes, the chlorite–sulfuric acid method is completely chlorine-free, although it suffers from the requirement of 25% more chlorite to produce an equivalent amount of chlorine dioxide. Alternatively, hydrogen peroxide may be efficiently used in small-scale applications.
Addition of sulfuric acid or any strong acid to chlorate salts produces chlorine dioxide.
Reduction of chlorate
In the laboratory, chlorine dioxide can also be prepared by reaction of potassium chlorate with oxalic acid:
or with oxalic and sulfuric acid:
Over 95% of the chlorine dioxide produced in the world today is made by reduction of sodium chlorate, for use in pulp bleaching. It is produced with high efficiency in a strong acid solution with a suitable reducing agent such as methanol, hydrogen peroxide, hydrochloric acid or sulfur dioxide. Modern technologies are based on methanol or hydrogen peroxide, as these chemistries allow the best economy and do not co-produce elemental chlorine. The overall reaction can be written as:
As a typical example, the reaction of sodium chlorate with hydrochloric acid in a single reactor is believed to proceed through the following pathway:
which gives the overall reaction
The commercially more important production route uses methanol as the reducing agent and sulfuric acid for the acidity. Two advantages of not using the chloride-based processes are that there is no formation of elemental chlorine, and that sodium sulfate, a valuable chemical for the pulp mill, is a side-product. These methanol-based processes provide high efficiency and can be made very safe.
The variant process using sodium chlorate, hydrogen peroxide and sulfuric acid has been increasingly used since 1999 for water treatment and other small-scale disinfection applications, since it produce a chlorine-free product at high efficiency, over 95%.
Other processes
Very pure chlorine dioxide can also be produced by electrolysis of a chlorite solution:
High-purity chlorine dioxide gas (7.7% in air or nitrogen) can be produced by the gas–solid method, which reacts dilute chlorine gas with solid sodium chlorite:
Handling properties
Chlorine dioxide is very different from elemental chlorine. One of the most important qualities of chlorine dioxide is its high water solubility, especially in cold water. Chlorine dioxide does not react with water; it remains a dissolved gas in solution. Chlorine dioxide is approximately 10 times more soluble in water than elemental chlorine but its solubility is very temperature-dependent.
At partial pressures above (or gas-phase concentrations greater than 10% volume in air at STP) of ClO2 may explosively decompose into chlorine and oxygen. The decomposition can be initiated by light, hot spots, chemical reaction, or pressure shock. Thus, chlorine dioxide is never handled as a pure gas, but is almost always handled in an aqueous solution in concentrations between 0.5 to 10 grams per liter. Its solubility increases at lower temperatures, so it is common to use chilled water (5 °C, 41 °F) when storing at concentrations above 3 grams per liter. In many countries, such as the United States, chlorine dioxide may not be transported at any concentration and is instead almost always produced on-site. In some countries, chlorine dioxide solutions below 3 grams per liter in concentration may be transported by land, but they are relatively unstable and deteriorate quickly.
Uses
Chlorine dioxide is used for bleaching of wood pulp and for the disinfection (called chlorination) of municipal drinking water, treatment of water in oil and gas applications, disinfection in the food industry, microbiological control in cooling towers, and textile bleaching. As a disinfectant, it is effective even at low concentrations because of its unique qualities.
Bleaching
Chlorine dioxide is sometimes used for bleaching of wood pulp in combination with chlorine, but it is used alone in ECF (elemental chlorine-free) bleaching sequences. It is used at moderately acidic pH (3.5 to 6). The use of chlorine dioxide minimizes the amount of organochlorine compounds produced. Chlorine dioxide (ECF technology) currently is the most important bleaching method worldwide. About 95% of all bleached kraft pulp is made using chlorine dioxide in ECF bleaching sequences.
Chlorine dioxide has been used to bleach flour.
Water treatment
The water treatment plant at Niagara Falls, New York first used chlorine dioxide for drinking water treatment in 1944 for destroying "taste and odor producing phenolic compounds." Chlorine dioxide was introduced as a drinking water disinfectant on a large scale in 1956, when Brussels, Belgium, changed from chlorine to chlorine dioxide. Its most common use in water treatment is as a pre-oxidant prior to chlorination of drinking water to destroy natural water impurities that would otherwise produce trihalomethanes upon exposure to free chlorine. Trihalomethanes are suspected carcinogenic disinfection by-products associated with chlorination of naturally occurring organics in raw water. Chlorine dioxide also produces 70% fewer halomethanes in the presence of natural organic matter compared to when elemental chlorine or bleach is used.
Chlorine dioxide is also superior to chlorine when operating above pH 7, in the presence of ammonia and amines, and for the control of biofilms in water distribution systems. Chlorine dioxide is used in many industrial water treatment applications as a biocide, including cooling towers, process water, and food processing.
Chlorine dioxide is less corrosive than chlorine and superior for the control of Legionella bacteria.
Chlorine dioxide is superior to some other secondary water disinfection methods, in that chlorine dioxide is not negatively impacted by pH, does not lose efficacy over time, because the bacteria will not grow resistant to it, and is not negatively impacted by silica and phosphates, which are commonly used potable water corrosion inhibitors. In the United States, it is an EPA-registered biocide.
It is more effective as a disinfectant than chlorine in most circumstances against waterborne pathogenic agents such as viruses, bacteria, and protozoa – including the cysts of Giardia and the oocysts of Cryptosporidium.
The use of chlorine dioxide in water treatment leads to the formation of the by-product chlorite, which is currently limited to a maximum of 1 part per million in drinking water in the USA. This EPA standard limits the use of chlorine dioxide in the US to relatively high-quality water, because this minimizes chlorite concentration, or water that is to be treated with iron-based coagulants, because iron can reduce chlorite to chloride. The World Health Organization also advises a 1ppm dosification.
Use in public crises
Chlorine dioxide has many applications as an oxidizer or disinfectant. Chlorine dioxide can be used for air disinfection and was the principal agent used in the decontamination of buildings in the United States after the 2001 anthrax attacks. After the disaster of Hurricane Katrina in New Orleans, Louisiana, and the surrounding Gulf Coast, chlorine dioxide was used to eradicate dangerous mold from houses inundated by the flood water.
In addressing the COVID-19 pandemic, the U.S. Environmental Protection Agency has posted a list of many disinfectants that meet its criteria for use in environmental measures against the causative coronavirus. Some are based on sodium chlorite that is activated into chlorine dioxide, though differing formulations are used in each product. Many other products on the EPA list contain sodium hypochlorite, which is similar in name but should not be confused with sodium chlorite because they have very different modes of chemical action.
Other disinfection uses
Chlorine dioxide may be used as a fumigant treatment to "sanitize" fruits such as blueberries, raspberries, and strawberries that develop molds and yeast.
Chlorine dioxide may be used to disinfect poultry by spraying or immersing it after slaughtering.
Chlorine dioxide may be used for the disinfection of endoscopes, such as under the trade name Tristel. It is also available in a trio consisting of a preceding pre-clean with surfactant and a succeeding rinse with deionized water and a low-level antioxidant.
Chlorine dioxide may be used for control of zebra and quagga mussels in water intakes.
Chlorine dioxide was shown to be effective in bedbug eradication.
For water purification during camping, disinfecting tablets containing chlorine dioxide are more effective against pathogens than those using household bleach, but typically cost more.
Other uses
Chlorine dioxide is used as an oxidant for destroying phenols in wastewater streams and for odor control in the air scrubbers of animal byproduct (rendering) plants. It is also available for use as a deodorant for cars and boats, in chlorine dioxide-generating packages that are activated by water and left in the boat or car overnight.
In dilute concentrations, chlorine dioxide is an ingredient that acts as an antiseptic agent in some mouthwashes.
Safety issues in water and supplements
Potential hazards with chlorine dioxide include poisoning and the risk of spontaneous ignition or explosion on contact with flammable materials.
Chlorine dioxide is toxic, and limits on human exposure are required to ensure its safe use. The United States Environmental Protection Agency has set a maximum level of 0.8 mg/L for chlorine dioxide in drinking water. The Occupational Safety and Health Administration (OSHA), an agency of the United States Department of Labor, has set an 8-hour permissible exposure limit of 0.1 ppm in air (0.3 mg/m3) for people working with chlorine dioxide.
Chlorine dioxide has been fraudulently and illegally marketed as an ingestible cure for a wide range of diseases, including childhood autism and coronavirus. Children who have been given enemas of chlorine dioxide as a supposed cure for childhood autism have suffered life-threatening ailments. The U.S. Food and Drug Administration (FDA) has stated that ingestion or other internal use of chlorine dioxide, outside of supervised oral rinsing using dilute concentrations, has no health benefits of any kind, and it should not be used internally for any reason.
Pseudomedicine
On 30 July and 1 October 2010, the United States Food and Drug Administration warned against the use of the product "Miracle Mineral Supplement", or "MMS", which when prepared according to the instructions produces chlorine dioxide. MMS has been marketed as a treatment for a variety of conditions, including HIV, cancer, autism, acne, and, more recently, COVID-19. Many have complained to the FDA, reporting life-threatening reactions, and even death. The FDA has warned consumers that MMS can cause serious harm to health, and stated that it has received numerous reports of nausea, diarrhea, severe vomiting, and life-threatening low blood pressure caused by dehydration. This warning was repeated for a third time on 12 August 2019, and a fourth on 8 April 2020, stating that ingesting MMS is just as hazardous as ingesting bleach, and urging consumers not to use them or give these products to their children for any reason, as there is no scientific evidence showing that chlorine dioxide has any beneficial medical properties.
| Physical sciences | Covalent oxides | Chemistry |
458680 | https://en.wikipedia.org/wiki/Myctophiformes | Myctophiformes | The Myctophiformes are an order of ray-finned fishes consisting of two families of deep-sea marine fish, most notably the highly abundant lanternfishes (Myctophidae). The blackchins (Neoscopelidae) contain six species in three genera, while the bulk of the family belongs to the Myctophidae, with over 30 genera and some 252 species.
The scientific name ultimately derives from Ancient Greek myktér (μυκτήρ, "nose") + óphis (ὄφῖς, "serpent") + Latin forma ("external form"), the Greek part in reference to the long, slender, and heavy-headed shape of these fishes.
Description and ecology
These smallish fishes inhabit the pelagic and benthopelagic zones of the deep sea. They are laterally compressed and usually have photophores (light organs). The eyes are large, in some decidedly huge, and generally directed straight sideways. The mouth also quite large and located at the tip of the snout; its gape extends to below the eyes or even beyond. They have an adipose fin. The pelvic fin has eight rays in most myctophiforms, and the number of branchiostegal rays is usually higher than six and lower than 12.
Systematics
The two families of the Myctophiformes are:
Myctophidae – lanternfishes
Neoscopelidae – blackchins
The extinct family Sardinioididae (containing Sardinioides and possibly Volcichthys) is known from the Late Cretaceous (Cenomanian to Campanian). The genus Neocassandra from the Late Paleocene is either considered its own family or a member of the Neoscopelidae. Some members were previously confused with aulopiforms.
The order Myctophiformes is anatomically similar to the grinners (Aulopiformes), but their pharyngobranchials and retractor muscles are more plesiomorphic. It was also allied with the more advanced spiny-rayed Teleostei (e.g. Paracanthopterygii) as "Ctenosquamata". These apomorphically have a fifth upper pharyngeal toothplate and a third internal levator muscle to move it, and molecular data also support the long-held view that these two lineages are at least closely related. Other sources ally them with the Lampriformes, which are often placed in a monotypic superorder "Lampridiomorpha". In a similar fashion, separation of the Myctophiformes in superorder "Scopelomorpha" has been proposed. The Aulopiformes, though, are usually considered to be closer or even among the Protacanthopterygii, one of the core groups of moderately advanced teleosts. As modern taxonomy tries to avoid a profusion of small taxa, and the delimitation of the Euteleostei (Protacanthopterygii sensu stricto and their allies) versus "Ctenosquamata" such as the Paracanthopterygii remains uncertain, the systematics and taxonomy of the Myctophiformes among the teleosts are in need of further study.
Timeline of genera
| Biology and health sciences | Myctophiformes | Animals |
458866 | https://en.wikipedia.org/wiki/Solid%20mechanics | Solid mechanics | Solid mechanics (also known as mechanics of solids) is the branch of continuum mechanics that studies the behavior of solid materials, especially their motion and deformation under the action of forces, temperature changes, phase changes, and other external or internal agents.
Solid mechanics is fundamental for civil, aerospace, nuclear, biomedical and mechanical engineering, for geology, and for many branches of physics and chemistry such as materials science. It has specific applications in many other areas, such as understanding the anatomy of living beings, and the design of dental prostheses and surgical implants. One of the most common practical applications of solid mechanics is the Euler–Bernoulli beam equation. Solid mechanics extensively uses tensors to describe stresses, strains, and the relationship between them.
Solid mechanics is a vast subject because of the wide range of solid materials available, such as steel, wood, concrete, biological materials, textiles, geological materials, and plastics.
Fundamental aspects
A solid is a material that can support a substantial amount of shearing force over a given time scale during a natural or industrial process or action. This is what distinguishes solids from fluids, because fluids also support normal forces which are those forces that are directed perpendicular to the material plane across from which they act and normal stress is the normal force per unit area of that material plane. Shearing forces in contrast with normal forces, act parallel rather than perpendicular to the material plane and the shearing force per unit area is called shear stress.
Therefore, solid mechanics examines the shear stress, deformation and the failure of solid materials and structures.
The most common topics covered in solid mechanics include:
stability of structures - examining whether structures can return to a given equilibrium after disturbance or partial/complete failure, see Structure mechanics
dynamical systems and chaos - dealing with mechanical systems highly sensitive to their given initial position
thermomechanics - analyzing materials with models derived from principles of thermodynamics
biomechanics - solid mechanics applied to biological materials e.g. bones, heart tissue
geomechanics - solid mechanics applied to geological materials e.g. ice, soil, rock
vibrations of solids and structures - examining vibration and wave propagation from vibrating particles and structures i.e. vital in mechanical, civil, mining, aeronautical, maritime/marine, aerospace engineering
fracture and damage mechanics - dealing with crack-growth mechanics in solid materials
composite materials - solid mechanics applied to materials made up of more than one compound e.g. reinforced plastics, reinforced concrete, fiber glass
variational formulations and computational mechanics - numerical solutions to mathematical equations arising from various branches of solid mechanics e.g. finite element method (FEM)
experimental mechanics - design and analysis of experimental methods to examine the behavior of solid materials and structures
Relationship to continuum mechanics
As shown in the following table, solid mechanics inhabits a central place within continuum mechanics. The field of rheology presents an overlap between solid and fluid mechanics.
Response models
A material has a rest shape and its shape departs away from the rest shape due to stress. The amount of departure from rest shape is called deformation, the proportion of deformation to original size is called strain. If the applied stress is sufficiently low (or the imposed strain is small enough), almost all solid materials behave in such a way that the strain is directly proportional to the stress; the coefficient of the proportion is called the modulus of elasticity. This region of deformation is known as the linearly elastic region.
It is most common for analysts in solid mechanics to use linear material models, due to ease of computation. However, real materials often exhibit non-linear behavior. As new materials are used and old ones are pushed to their limits, non-linear material models are becoming more common.
These are basic models that describe how a solid responds to an applied stress:
Elasticity – When an applied stress is removed, the material returns to its undeformed state. Linearly elastic materials, those that deform proportionally to the applied load, can be described by the linear elasticity equations such as Hooke's law.
Viscoelasticity – These are materials that behave elastically, but also have damping: when the stress is applied and removed, work has to be done against the damping effects and is converted in heat within the material resulting in a hysteresis loop in the stress–strain curve. This implies that the material response has time-dependence.
Plasticity – Materials that behave elastically generally do so when the applied stress is less than a yield value. When the stress is greater than the yield stress, the material behaves plastically and does not return to its previous state. That is, deformation that occurs after yield is permanent.
Viscoplasticity - Combines theories of viscoelasticity and plasticity and applies to materials like gels and mud.
Thermoelasticity - There is coupling of mechanical with thermal responses. In general, thermoelasticity is concerned with elastic solids under conditions that are neither isothermal nor adiabatic. The simplest theory involves the Fourier's law of heat conduction, as opposed to advanced theories with physically more realistic models.
Timeline
1452–1519 Leonardo da Vinci made many contributions
1638: Galileo Galilei published the book "Two New Sciences" in which he examined the failure of simple structures
1660: Hooke's law by Robert Hooke
1687: Isaac Newton published "Philosophiae Naturalis Principia Mathematica" which contains Newton's laws of motion
1750: Euler–Bernoulli beam equation
1700–1782: Daniel Bernoulli introduced the principle of virtual work
1707–1783: Leonhard Euler developed the theory of buckling of columns
1826: Claude-Louis Navier published a treatise on the elastic behaviors of structures
1873: Carlo Alberto Castigliano presented his dissertation "Intorno ai sistemi elastici", which contains his theorem for computing displacement as partial derivative of the strain energy. This theorem includes the method of least work as a special case
1874: Otto Mohr formalized the idea of a statically indeterminate structure.
1922: Timoshenko corrects the Euler–Bernoulli beam equation
1936: Hardy Cross' publication of the moment distribution method, an important innovation in the design of continuous frames.
1941: Alexander Hrennikoff solved the discretization of plane elasticity problems using a lattice framework
1942: R. Courant divided a domain into finite subregions
1956: J. Turner, R. W. Clough, H. C. Martin, and L. J. Topp's paper on the "Stiffness and Deflection of Complex Structures" introduces the name "finite-element method" and is widely recognized as the first comprehensive treatment of the method as it is known today
| Physical sciences | Solid mechanics | null |
2549117 | https://en.wikipedia.org/wiki/Variable%20capacitor | Variable capacitor | A variable capacitor is a capacitor whose capacitance may be intentionally and repeatedly changed mechanically or electronically. Variable capacitors are often used in L/C circuits to set the resonance frequency, e.g. to tune a radio (therefore it is sometimes called a tuning capacitor or tuning condenser), or as a variable reactance, e.g. for impedance matching in antenna tuners.
Mechanically controlled capacitance
In mechanically controlled variable capacitors, the distance between the plates, or the amount of plate surface area which overlaps, can be changed.
The most common form arranges a group of semicircular metal plates on a rotary axis ("rotor") that are positioned in the gaps between a set of stationary plates ("stator") so that the area of overlap can be changed by rotating the axis. Air or plastic foils can be used as dielectric material.
By choosing the shape of the rotary plates, various functions of capacitance vs. angle can be created, e.g. to obtain a linear frequency scale. Various forms of reduction gear mechanisms are often used to achieve finer tuning control, i.e. to spread the variation of capacity over a larger angle, often several turns. Maximum capacitance is achieved when the plates are "meshed" together, that is, they are inter-laced. Minimum capacitance is achieved when the plates are "unmeshed", that is, they are not inter-laced.
A vacuum variable capacitor uses a set of plates made from concentric cylinders that can be slid in or out of an opposing set of cylinders (sleeve and plunger). These plates are then sealed inside of a non-conductive envelope such as glass or ceramic and placed under a high vacuum. The movable part (plunger) is mounted on a flexible metal membrane that seals and maintains the vacuum. A screw shaft is attached to the plunger; when the shaft is turned the plunger moves in or out of the sleeve and the value of the capacitor changes. The vacuum not only increases the working voltage and current handling capacity of the capacitor, it also greatly reduces the chance of arcing across the plates. The most common usage for vacuum variables are in high-powered transmitters such as those used for broadcasting, military and amateur radio, as well as high-powered RF tuning networks. Vacuum variables can also be more convenient; since the elements are under a vacuum, the working voltage can be higher than an air variable the same size, allowing the size of the vacuum capacitor to be reduced.
Very cheap variable capacitors are constructed from layered aluminium and plastic foils that are variably pressed together using a screw. These so-called squeezers cannot provide a stable and reproducible capacitance, however. A variant of this structure that allows for linear movement of one set of plates to change the plate overlap area is also used and might be called a slider. This has practical advantages for makeshift or home construction, and may be found in resonant-loop antennas or crystal radios.
Small variable capacitors operated by screwdriver (for instance, to precisely set a resonant frequency at the factory and then never be adjusted again) are called trimmer capacitors. In addition to air and plastic, trimmers can also be made using a solid dielectric, such as mica.
Special forms of mechanically variable capacitors
Multiple sections
Very often, multiple stator/rotor sections are arranged behind one another on the same axis, allowing for several tuned circuits to be adjusted using the same control, e.g. a preselector, an input filter and the corresponding oscillator in a receiver circuit. The sections can have identical or different nominal capacitances, e.g. 2 × 330 pF for AM filter and oscillator, plus 3 × 45 pF for two filters and an oscillator in the FM section of the same receiver. Capacitors with multiple sections often include trimmer capacitors in parallel to the variable sections, used to adjust all tuned circuits to the same frequency.
Butterfly
A butterfly capacitor is a form of rotary variable capacitor with two independent sets of stator plates opposing each other, and a butterfly-shaped rotor arranged so that turning the rotor will vary the capacitances between the rotor and either stator equally.
Butterfly capacitors are used in symmetrical tuned circuits, e.g. RF power amplifier stages in push-pull configuration or symmetrical antenna tuners where the rotor needs to be "cold", i.e. connected to RF (but not necessarily DC) ground potential. Since the peak RF current normally flows from one stator to the other without going through wiper contacts, butterfly capacitors can handle large resonance RF currents, e.g. in magnetic loop antennas.
In a butterfly capacitor, the stators and each half of the rotor can only cover a maximum angle of 90° since there must be a position without rotor/stator overlap corresponding to minimum capacity, therefore a turn of only 90° covers the entire capacitance range.
Split stator
The closely related split stator variable capacitor does not have the limitation of 90° angle since it uses two separate packs of rotor electrodes arranged axially behind one another. Unlike in a capacitor with several sections, the rotor plates in a split stator capacitor are mounted on opposite sides of the rotor axis. While the split stator capacitor benefits from larger electrodes compared to the butterfly capacitor, as well as a rotation angle of up to 180°, the separation of rotor plates incurs some losses since RF current has to pass the rotor axis instead of flowing straight through each rotor vane.
Differential
Differential variable capacitors also have two independent stators, but unlike in the butterfly capacitor where capacities on both sides increase equally as the rotor is turned, in a differential variable capacitor one section's capacity will increase while the other section's decreases, keeping the sum of the two stator capacitances constant. Differential variable capacitors can therefore be used in capacitive potentiometric circuits.
History
The variable capacitor with air dielectric was invented by the Hungarian engineer Dezső Korda. He received a German patent for the invention on 13 December 1893.
Electronically controlled capacitance
Voltage tuned capacitance
The thickness of the depletion layer of a reverse-biased semiconductor diode varies with the DC voltage applied across the diode. Any diode exhibits this effect (including p/n junctions in transistors), but devices specifically sold as variable capacitance diodes (also called varactors or varicaps) are designed with a large junction area and a doping profile specifically designed to maximize capacitance.
Their use is limited to low signal amplitudes to avoid obvious distortions as the capacitance would be affected by the change of signal voltage, precluding their use in the input stages of high-quality RF communications receivers, where they would add unacceptable levels of intermodulation. At VHF/UHF frequencies, e.g. in FM Radio or TV tuners, dynamic range is limited by noise rather than large signal handling requirements, and varicaps are commonly used in the signal path.
Varicaps are used for frequency modulation of oscillators, and to make high-frequency voltage controlled oscillators (VCOs), the core component in phase-locked loop (PLL) frequency synthesizers that are ubiquitous in modern communications equipment.
BST device are based on Barium Strontium Titanate and vary the capacitance by applying high voltage to the device. They have a dedicated analogue control input and therefore introduce less non-linearities than varactor diodes, especially for higher signal voltages. The limitations for BST are stability over temperature and linearity in demanding applications.
Digitally tuned capacitance
A digitally tuned capacitor is an IC variable capacitor based on several technologies. MEMS, BST and SOI/SOS devices are available from a number of suppliers and vary in capacitance range, quality factor and resolution for different RF tuning applications.
MEMS devices have the highest quality factor and are highly linear, and therefore are suitable for antenna aperture tuning, dynamic impedance matching, power amplifier load matching and adjustable filters. RF tuning MEMS are still a relatively new technology and has not yet been accepted broadly.
SOI/SOS tuning devices are constructed as solid state FET switches built on insulated CMOS wafers and use MIM caps arranged in binary-weighted values to achieve different capacitance values. SOI/SOS switches have high linearity and are well suited to low power applications where high voltages are not present. High voltage endurance requires multiple FET devices in series which adds series resistance and lowers the quality factor.
The capacitance values are designed for antenna impedance matching in multi-band LTE GSM/WCDMA cellular handsets and mobile TV receivers that operate over wide frequency ranges, such as the European DVB-H and Japanese ISDB-T mobile TV systems.
Transducers
Variable capacitance is sometimes used to convert physical phenomena into electrical signals.
In a capacitor microphone (commonly known as a condenser microphone), the diaphragm acts as one plate of a capacitor, and vibrations produce changes in the distance between the diaphragm and a fixed plate, changing the voltage maintained across the capacitor plates.
Some types of industrial sensors use a capacitor element to convert physical quantities such as pressure, displacement or relative humidity to an electrical signal for measurement purposes.
Capacitive sensors can also be used in the place of switches, e.g. in computer keyboards or "touch buttons" for elevators that have no user-movable parts.
| Technology | Components | null |
2549543 | https://en.wikipedia.org/wiki/Thermal%20power%20station | Thermal power station | A thermal power station, also known as a thermal power plant, is a type of power station in which the heat energy generated from various fuel sources (e.g., coal, natural gas, nuclear fuel, etc.) is converted to electrical energy. The heat from the source is converted into mechanical energy using a thermodynamic power cycle (such as a Diesel cycle, Rankine cycle, Brayton cycle, etc.). The most common cycle involves a working fluid (often water) heated and boiled under high pressure in a pressure vessel to produce high-pressure steam. This high pressure-steam is then directed to a turbine, where it rotates the turbine's blades. The rotating turbine is mechanically connected to an electric generator which converts rotary motion into electricity. Fuels such as natural gas or oil can also be burnt directly in gas turbines (internal combustion), skipping the steam generation step. These plants can be of the open cycle or the more efficient combined cycle type.
The majority of the world's thermal power stations are driven by steam turbines, gas turbines, or a combination of the two. The efficiency of a thermal power station is determined by how effectively it converts heat energy into electrical energy, specifically the ratio of saleable electricity to the heating value of the fuel used. Different thermodynamic cycles have varying efficiencies, with the Rankine cycle generally being more efficient than the Otto or Diesel cycles. In the Rankine cycle, the low-pressure exhaust from the turbine enters a steam condenser where it is cooled to produce hot condensate which is recycled to the heating process to generate even more high pressure steam.
The design of thermal power stations depends on the intended energy source. In addition to fossil and nuclear fuel, some stations use geothermal power, solar energy, biofuels, and waste incineration. Certain thermal power stations are also designed to produce heat for industrial purposes, provide district heating, or desalinate water, in addition to generating electrical power. Emerging technologies such as supercritical and ultra-supercritical thermal power stations operate at higher temperatures and pressures for increased efficiency and reduced emissions. Cogeneration or CHP (Combined Heat and Power) technology, the simultaneous production of electricity and useful heat from the same fuel source, improves the overall efficiency by using waste heat for heating purposes. Older, less efficient thermal power stations are being decommissioned or adapted to use cleaner and renewable energy sources.
Thermal power stations produce 70% of the world's electricity. They often provide reliable, stable, and continuous baseload power supply essential for economic growth. They ensure energy security by maintaining grid stability, especially in regions where they complement intermittent renewable energy sources dependent on weather conditions. The operation of thermal power stations contributes to the local economy by creating jobs in construction, maintenance, and fuel extraction industries. On the other hand, burning of fossil fuels releases greenhouse gases (contributing to climate change) and air pollutants such as sulfur oxides and nitrogen oxides (leading to acid rain and respiratory diseases). Carbon capture and storage (CCS) technology can reduce the greenhouse gas emissions of fossil-fuel-based thermal power stations, however it is expensive and has seldom been implemented. Government regulations and international agreements are being enforced to reduce harmful emissions and promote cleaner power generation.
Types of thermal energy
Almost all coal-fired power stations, petroleum, nuclear, geothermal, solar thermal electric, and waste incineration plants, as well as all natural gas power stations are thermal. Natural gas is frequently burned in gas turbines as well as boilers. The waste heat from a gas turbine, in the form of hot exhaust gas, can be used to raise steam by passing this gas through a heat recovery steam generator (HRSG). The steam is then used to drive a steam turbine in a combined cycle plant that improves overall efficiency. Power stations burning coal, fuel oil, or natural gas are often called fossil fuel power stations. Some biomass-fueled thermal power stations have appeared also. Non-nuclear thermal power stations, particularly fossil-fueled plants, which do not use cogeneration are sometimes referred to as conventional power stations.
Commercial electric utility power stations are usually constructed on a large scale and designed for continuous operation. Virtually all electric power stations use three-phase electrical generators to produce alternating current (AC) electric power at a frequency of 50 Hz or 60 Hz. Large companies or institutions may have their own power stations to supply heating or electricity to their facilities, especially if steam is created anyway for other purposes. Steam-driven power stations have been used to drive most ships in most of the 20th century. Shipboard power stations usually directly couple the turbine to the ship's propellers through gearboxes. Power stations in such ships also provide steam to smaller turbines driving electric generators to supply electricity. Nuclear marine propulsion is, with few exceptions, used only in naval vessels. There have been many turbo-electric ships in which a steam-driven turbine drives an electric generator which powers an electric motor for propulsion.
Cogeneration plants, often called combined heat and power (CHP) facilities, produce both electric power and heat for process heat or space heating, such as steam and hot water.
History
The reciprocating steam engine has been used to produce mechanical power since the 18th century, with notable improvements being made by James Watt. When the first commercially developed central electrical power stations were established in 1882 at Pearl Street Station in New York and Holborn Viaduct power station in London, reciprocating steam engines were used. The development of the steam turbine in 1884 provided larger and more efficient machine designs for central generating stations. By 1892 the turbine was considered a better alternative to reciprocating engines; turbines offered higher speeds, more compact machinery, and stable speed regulation allowing for parallel synchronous operation of generators on a common bus. After about 1905, turbines entirely replaced reciprocating engines in almost all large central power stations.
The largest reciprocating engine-generator sets ever built were completed in 1901 for the Manhattan Elevated Railway. Each of seventeen units weighed about 500 tons and was rated 6000 kilowatts; a contemporary turbine set of similar rating would have weighed about 20% as much.
Thermal power generation efficiency
The energy efficiency of a conventional thermal power station is defined as saleable energy produced as a percent of the heating value of the fuel consumed. A simple cycle gas turbine achieves energy conversion efficiencies from 20 to 35%. Typical coal-based power plants operating at steam pressures of 170 bar and 570 °C run at efficiency of 35 to 38%, with state-of-the-art fossil fuel plants at 46% efficiency. Combined-cycle systems can reach higher values. As with all heat engines, their efficiency is limited, and governed by the laws of thermodynamics.
The Carnot efficiency dictates that higher efficiencies can be attained by increasing the temperature of the steam. Sub-critical pressure fossil fuel power stations can achieve 36–40% efficiency. Supercritical designs have efficiencies in the low to mid 40% range, with new "ultra critical" designs using pressures above and multiple stage reheat reaching 45–48% efficiency. Above the critical point for water of and , there is no phase transition from water to steam, but only a gradual decrease in density.
Currently most nuclear power stations must operate below the temperatures and pressures that coal-fired plants do, in order to provide more conservative safety margins within the systems that remove heat from the nuclear fuel. This, in turn, limits their thermodynamic efficiency to 30–32%. Some advanced reactor designs being studied, such as the very-high-temperature reactor, Advanced Gas-cooled Reactor, and supercritical water reactor, would operate at temperatures and pressures similar to current coal plants, producing comparable thermodynamic efficiency.
The energy of a thermal power station not utilized in power production must leave the plant in the form of heat to the environment. This waste heat can go through a condenser and be disposed of with cooling water or in cooling towers. If the waste heat is instead used for district heating, it is called cogeneration. An important class of thermal power station is that associated with desalination facilities; these are typically found in desert countries with large supplies of natural gas, and in these plants freshwater production and electricity are equally important co-products.
Other types of power stations are subject to different efficiency limitations. Most hydropower stations in the United States are about 90 percent efficient in converting the energy of falling water into electricity while the efficiency of a wind turbine is limited by Betz's law, to about 59.3%, and actual wind turbines show lower efficiency.
Electricity cost
The direct cost of electric energy produced by a thermal power station is the result of cost of fuel, capital cost for the plant, operator labour, maintenance, and such factors as ash handling and disposal. Indirect social or environmental costs, such as the economic value of environmental impacts, or environmental and health effects of the complete fuel cycle and plant decommissioning, are not usually assigned to generation costs for thermal stations in utility practice, but may form part of an environmental impact assessment. Those indirect costs belong to the broader concept of externalities.
Boiler and steam cycle
In the nuclear plant field, steam generator refers to a specific type of large heat exchanger used in a pressurized water reactor (PWR) to thermally connect the primary (reactor plant) and secondary (steam plant) systems, which generates steam. In a boiling water reactor (BWR), no separate steam generator is used and water boils in the reactor core.
In some industrial settings, there can also be steam-producing heat exchangers called heat recovery steam generators (HRSG) which utilize heat from some industrial process, most commonly utilizing hot exhaust from a gas turbine. The steam generating boiler has to produce steam at the high purity, pressure and temperature required for the steam turbine that drives the electrical generator.
Geothermal plants do not need boilers because they use naturally occurring steam sources. Heat exchangers may be used where the geothermal steam is very corrosive or contains excessive suspended solids.
A fossil fuel steam generator includes an economizer, a steam drum, and the furnace with its steam generating tubes and superheater coils. Necessary safety valves are located at suitable points to protect against excessive boiler pressure. The air and flue gas path equipment include: forced draft (FD) fan, air preheater (AP), boiler furnace, induced draft (ID) fan, fly ash collectors (electrostatic precipitator or baghouse), and the flue-gas stack.
Feed water heating
The boiler feed water used in the steam boiler is a means of transferring heat energy from the burning fuel to the mechanical energy of the spinning steam turbine. The total feed water consists of recirculated condensate water and purified makeup water. Because the metallic materials it contacts are subject to corrosion at high temperatures and pressures, the makeup water is highly purified before use. A system of water softeners and ion exchange demineralizes produces water so pure that it coincidentally becomes an electrical insulator, with conductivity in the range of 0.3–1.0 microsiemens per centimeter. The makeup water in a 500 MWe plant amounts to perhaps 120 US gallons per minute (7.6 L/s) to replace water drawn off from the boiler drums for water purity management, and to also offset the small losses from steam leaks in the system.
The feed water cycle begins with condensate water being pumped out of the condenser after traveling through the steam turbines. The condensate flow rate at full load in a 500 MW plant is about 6,000 US gallons per minute (400 L/s).
The water is usually pressurized in two stages, and typically flows through a series of six or seven intermediate feed water heaters, heated up at each point with steam extracted from an appropriate extraction connection on the turbines and gaining temperature at each stage. Typically, in the middle of this series of feedwater heaters, and before the second stage of pressurization, the condensate plus the makeup water flows through a deaerator that removes dissolved air from the water, further purifying and reducing its corrosiveness. The water may be dosed following this point with hydrazine, a chemical that removes the remaining oxygen in the water to below 5 parts per billion (ppb). It is also dosed with pH control agents such as ammonia or morpholine to keep the residual acidity low and thus non-corrosive.
Boiler operation
The boiler is a rectangular furnace about on a side and tall. Its walls are made of a web of high pressure steel tubes about in diameter.
Fuel such as pulverized coal is air-blown into the furnace through burners located at the four corners, or along one wall, or two opposite walls, and it is ignited to rapidly burn, forming a large fireball at the center. The thermal radiation of the fireball heats the water that circulates through the boiler tubes near the boiler perimeter. The water circulation rate in the boiler is three to four times the throughput. As the water in the boiler circulates it absorbs heat and changes into steam. It is separated from the water inside a drum at the top of the furnace. The saturated steam is introduced into superheat pendant tubes that hang in the hottest part of the combustion gases as they exit the furnace. Here the steam is superheated to to prepare it for the turbine.
Plants that use gas turbines to heat the water for conversion into steam use boilers known as heat recovery steam generators (HRSG). The exhaust heat from the gas turbines is used to make superheated steam that is then used in a conventional water-steam generation cycle, as described in the gas turbine combined-cycle plants section.
Boiler furnace and steam drum
The water enters the boiler through a section in the convection pass called the economizer. From the economizer it passes to the steam drum and from there it goes through downcomers to inlet headers at the bottom of the water walls. From these headers the water rises through the water walls of the furnace where some of it is turned into steam and the mixture of water and steam then re-enters the steam drum. This process may be driven purely by natural circulation (because the water is the downcomers is denser than the water/steam mixture in the water walls) or assisted by pumps. In the steam drum, the water is returned to the downcomers and the steam is passed through a series of steam separators and dryers that remove water droplets from the steam. The dry steam then flows into the superheater coils.
The boiler furnace auxiliary equipment includes coal feed nozzles and igniter guns, soot blowers, water lancing, and observation ports (in the furnace walls) for observation of the furnace interior. Furnace explosions due to any accumulation of combustible gases after a trip-out are avoided by flushing out such gases from the combustion zone before igniting the coal.
The steam drum (as well as the superheater coils and headers) have air vents and drains needed for initial start up.
Superheater
Fossil fuel power stations often have a superheater section in the steam generating furnace. The steam passes through drying equipment inside the steam drum on to the superheater, a set of tubes in the furnace. Here the steam picks up more energy from hot flue gases outside the tubing, and its temperature is now superheated above the saturation temperature. The superheated steam is then piped through the main steam lines to the valves before the high-pressure turbine.
Nuclear-powered steam plants do not have such sections but produce steam at essentially saturated conditions. Experimental nuclear plants were equipped with fossil-fired superheaters in an attempt to improve overall plant operating cost.
Steam condensing
The condenser condenses the steam from the exhaust of the turbine into liquid to allow it to be pumped. If the condenser can be made cooler, the pressure of the exhaust steam is reduced and efficiency of the cycle increases.
The surface condenser is a shell and tube heat exchanger in which cooling water is circulated through the tubes. The exhaust steam from the low-pressure turbine enters the shell, where it is cooled and converted to condensate (water) by flowing over the tubes as shown in the adjacent diagram. Such condensers use steam ejectors or rotary motor-driven exhausts for continuous removal of air and gases from the steam side to maintain vacuum.
For best efficiency, the temperature in the condenser must be kept as low as practical in order to achieve the lowest possible pressure in the condensing steam. Since the condenser temperature can almost always be kept significantly below 100 °C where the vapor pressure of water is much less than atmospheric pressure, the condenser generally works under vacuum. Thus leaks of non-condensible air into the closed loop must be prevented.
Typically the cooling water causes the steam to condense at a temperature of about and that creates an absolute pressure in the condenser of about , i.e. a vacuum of about relative to atmospheric pressure. The large decrease in volume that occurs when water vapor condenses to liquid creates the vacuum that generally increases the efficiency of the turbines.
The limiting factor is the temperature of the cooling water and that, in turn, is limited by the prevailing average climatic conditions at the power station's location (it may be possible to lower the temperature beyond the turbine limits during winter, causing excessive condensation in the turbine). Plants operating in hot climates may have to reduce output if their source of condenser cooling water becomes warmer; unfortunately this usually coincides with periods of high electrical demand for air conditioning.
The condenser generally uses either circulating cooling water from a cooling tower to reject waste heat to the atmosphere, or once-through cooling (OTC) water from a river, lake or ocean. In the United States, about two-thirds of power plants use OTC systems, which often have significant adverse environmental impacts. The impacts include thermal pollution and killing large numbers of fish and other aquatic species at cooling water intakes.
The heat absorbed by the circulating cooling water in the condenser tubes must also be removed to maintain the ability of the water to cool as it circulates. This is done by pumping the warm water from the condenser through either natural draft, forced draft or induced draft cooling towers (as seen in the adjacent image) that reduce the temperature of the water by evaporation, by about —expelling waste heat to the atmosphere. The circulation flow rate of the cooling water in a 500 MW unit is about 14.2 m3/s (500 ft3/s or 225,000 US gal/min) at full load.
The condenser tubes are typically made stainless steel or other alloys to resist corrosion from either side. Nevertheless, they may become internally fouled during operation by bacteria or algae in the cooling water or by mineral scaling, all of which inhibit heat transfer and reduce thermodynamic efficiency. Many plants include an automatic cleaning system that circulates sponge rubber balls through the tubes to scrub them clean without the need to take the system off-line.
The cooling water used to condense the steam in the condenser returns to its source without having been changed other than having been warmed. If the water returns to a local water body (rather than a circulating cooling tower), it is often tempered with cool 'raw' water to prevent thermal shock when discharged into that body of water.
Another form of condensing system is the air-cooled condenser. The process is similar to that of a radiator and fan. Exhaust heat from the low-pressure section of a steam turbine runs through the condensing tubes, the tubes are usually finned and ambient air is pushed through the fins with the help of a large fan. The steam condenses to water to be reused in the water-steam cycle. Air-cooled condensers typically operate at a higher temperature than water-cooled versions. While saving water, the efficiency of the cycle is reduced (resulting in more carbon dioxide per megawatt-hour of electricity).
From the bottom of the condenser, powerful condensate pumps recycle the condensed steam (water) back to the water/steam cycle.
Reheater
Power station furnaces may have a reheater section containing tubes heated by hot flue gases outside the tubes. Exhaust steam from the high-pressure turbine is passed through these heated tubes to collect more energy before driving the intermediate and then low-pressure turbines.
Air path
External fans are provided to give sufficient air for combustion. The Primary air fan takes air from the atmosphere and, first warms the air in the air preheater for better economy. Primary air then passes through the coal pulverizers, and carries the coal dust to the burners for injection into the furnace. The Secondary air fan takes air from the atmosphere and, first warms the air in the air preheater for better economy. Secondary air is mixed with the coal/primary air flow in the burners.
The induced draft fan assists the FD fan by drawing out combustible gases from the furnace, maintaining slightly below atmospheric pressure in the furnace to avoid leakage of combustion products from the boiler casing.
Steam turbine generator
A steam turbine generator consists of a series of steam turbines interconnected to each other and a generator on a common shaft.
Steam turbine
There is usually a high-pressure turbine at one end, followed by an intermediate-pressure turbine, and finally one, two, or three low-pressure turbines, and the shaft that connects to the generator. As steam moves through the system and loses pressure and thermal energy, it expands in volume, requiring increasing diameter and longer blades at each succeeding stage to extract the remaining energy. The entire rotating mass may be over 200 metric tons and long. It is so heavy that it must be kept turning slowly even when shut down (at 3 rpm) so that the shaft will not bow even slightly and become unbalanced. This is so important that it is one of only six functions of blackout emergency power batteries on site. (The other five being emergency lighting, communication, station alarms, generator hydrogen seal system, and turbogenerator lube oil.)
For a typical late 20th-century power station, superheated steam from the boiler is delivered through piping at and to the high-pressure turbine, where it falls in pressure to and to in temperature through the stage. It exits via cold reheat lines and passes back into the boiler, where the steam is reheated in special reheat pendant tubes back to . The hot reheat steam is conducted to the intermediate-pressure turbine, where it falls in both temperature and pressure and exits directly to the long-bladed low-pressure turbines and finally exits to the condenser.
Turbo generator
The generator, typically about long and in diameter, contains a stationary stator and a spinning rotor, each containing miles of heavy copper conductor. There is generally no permanent magnet, thus preventing black starts. In operation it generates up to 21,000 amperes at 24,000 volts AC (504 MWe) as it spins at either 3,000 or 3,600 rpm, synchronized to the power grid. The rotor spins in a sealed chamber cooled with hydrogen gas, selected because it has the highest known heat transfer coefficient of any gas and for its low viscosity, which reduces windage losses. This system requires special handling during startup, with air in the chamber first displaced by carbon dioxide before filling with hydrogen. This ensures that a highly explosive hydrogen–oxygen environment is not created.
The power grid frequency is 60 Hz across North America and 50 Hz in Europe, Oceania, Asia (Korea and parts of Japan are notable exceptions), and parts of Africa. The desired frequency affects the design of large turbines, since they are highly optimized for one particular speed.
The electricity flows to a distribution yard where transformers increase the voltage for transmission to its destination.
The steam turbine-driven generators have auxiliary systems enabling them to work satisfactorily and safely. The steam turbine generator, being rotating equipment, generally has a heavy, large-diameter shaft. The shaft therefore requires not only supports but also has to be kept in position while running. To minimize the frictional resistance to the rotation, the shaft has a number of bearings. The bearing shells, in which the shaft rotates, are lined with a low-friction material like Babbitt metal. Oil lubrication is provided to further reduce the friction between shaft and bearing surface and to limit the heat generated.
Stack gas path and cleanup
As the combustion flue gas exits the boiler it is routed through a rotating flat basket of metal mesh which picks up heat and returns it to incoming fresh air as the basket rotates. This is called the air preheater. The gas exiting the boiler is laden with fly ash, which are tiny spherical ash particles. The flue gas contains nitrogen along with combustion products carbon dioxide, sulfur dioxide, and nitrogen oxides. The fly ash is removed by fabric bag filters in baghouses or electrostatic precipitators. Once removed, the fly ash byproduct can sometimes be used in the manufacturing of concrete. This cleaning up of flue gases, however, only occurs in plants that are fitted with the appropriate technology. Still, the majority of coal-fired power stations in the world do not have these facilities. Legislation in Europe has been efficient to reduce flue gas pollution. Japan has been using flue gas cleaning technology for over 30 years and the US has been doing the same for over 25 years. China is now beginning to grapple with the pollution caused by coal-fired power stations.
Where required by law, the sulfur and nitrogen oxide pollutants are removed by stack gas scrubbers which use a pulverized limestone or other alkaline wet slurry to remove those pollutants from the exit stack gas. Other devices use catalysts to remove nitrous oxide compounds from the flue-gas stream. The gas travelling up the flue-gas stack may by this time have dropped to about . A typical flue-gas stack may be tall to disperse the remaining flue gas components in the atmosphere. The tallest flue-gas stack in the world is tall at the Ekibastuz GRES-2 Power Station in Kazakhstan.
In the United States and a number of other countries, atmospheric dispersion modeling studies are required to determine the flue-gas stack height needed to comply with the local air pollution regulations. The United States also requires the height of a flue-gas stack to comply with what is known as the "good engineering practice" (GEP) stack height. In the case of existing flue gas stacks that exceed the GEP stack height, any air pollution dispersion modeling studies for such stacks must use the GEP stack height rather than the actual stack height.
Carbon capture and storage (CCS) captures carbon dioxide from the flue gas of power plants or other industry, transporting it to an appropriate location where it can be buried securely in an underground reservoir. Between 1972 and 2017, plans were made to add CCS to enough coal and gas power plants to sequester 171 million tonnes of per year, but by 2021 over 98% of these plans had failed. Cost, the absence of measures to address long-term liability for stored CO2, and limited social acceptability have all contributed to project cancellations. As of 2024, CCS is in operation at only five power plants worldwide.
Auxiliary systems
Boiler make-up water treatment plant and storage
Since there is continuous withdrawal of steam and continuous return of condensate to the boiler, losses due to blowdown and leakages have to be made up to maintain a desired water level in the boiler steam drum. For this, continuous make-up water is added to the boiler water system. Impurities in the raw water input to the plant generally consist of calcium and magnesium salts which impart hardness to the water. Hardness in the make-up water to the boiler will form deposits on the tube water surfaces which will lead to overheating and failure of the tubes. Thus, the salts have to be removed from the water, and that is done by a water demineralising treatment plant (DM). A DM plant generally consists of cation, anion, and mixed bed exchangers. Any ions in the final water from this process consist essentially of hydrogen ions and hydroxide ions, which recombine to form pure water. Very pure DM water becomes highly corrosive once it absorbs oxygen from the atmosphere because of its very high affinity for oxygen.
The capacity of the DM plant is dictated by the type and quantity of salts in the raw water input. However, some storage is essential as the DM plant may be down for maintenance. For this purpose, a storage tank is installed from which DM water is continuously withdrawn for boiler make-up. The storage tank for DM water is made from materials not affected by corrosive water, such as PVC. The piping and valves are generally of stainless steel. Sometimes, a steam blanketing arrangement or stainless steel doughnut float is provided on top of the water in the tank to avoid contact with air. DM water make-up is generally added at the steam space of the surface condenser (i.e., the vacuum side). This arrangement not only sprays the water but also DM water gets deaerated, with the dissolved gases being removed by a de-aerator through an ejector attached to the condenser.
Fuel preparation system
In coal-fired power stations, the raw feed coal from the coal storage area is first crushed into small pieces and then conveyed to the coal feed hoppers at the boilers. The coal is next pulverized into a very fine powder. The pulverizers may be ball mills, rotating drum grinders, or other types of grinders.
Some power stations burn fuel oil rather than coal. The oil must kept warm (above its pour point) in the fuel oil storage tanks to prevent the oil from congealing and becoming unpumpable. The oil is usually heated to about 100 °C before being pumped through the furnace fuel oil spray nozzles.
Boilers in some power stations use processed natural gas as their main fuel. Other power stations may use processed natural gas as auxiliary fuel in the event that their main fuel supply (coal or oil) is interrupted. In such cases, separate gas burners are provided on the boiler furnaces.
Barring gear
Barring gear (or "turning gear") is the mechanism provided to rotate the turbine generator shaft at a very low speed after unit stoppages. Once the unit is "tripped" (i.e., the steam inlet valve is closed), the turbine coasts down towards standstill. When it stops completely, there is a tendency for the turbine shaft to deflect or bend if allowed to remain in one position too long. This is because the heat inside the turbine casing tends to concentrate in the top half of the casing, making the top half portion of the shaft hotter than the bottom half. The shaft therefore could warp or bend by millionths of inches.
This small shaft deflection, only detectable by eccentricity meters, would be enough to cause damaging vibrations to the entire steam turbine generator unit when it is restarted. The shaft is therefore automatically turned at low speed (about one percent rated speed) by the barring gear until it has cooled sufficiently to permit a complete stop.
Oil system
An auxiliary oil system pump is used to supply oil at the start-up of the steam turbine generator. It supplies the hydraulic oil system required for steam turbine's main inlet steam stop valve, the governing control valves, the bearing and seal oil systems, the relevant hydraulic relays and other mechanisms.
At a preset speed of the turbine during start-ups, a pump driven by the turbine main shaft takes over the functions of the auxiliary system.
Generator cooling
While small generators may be cooled by air drawn through filters at the inlet, larger units generally require special cooling arrangements. Hydrogen gas cooling, in an oil-sealed casing, is used because it has the highest known heat transfer coefficient of any gas and for its low viscosity which reduces windage losses. This system requires special handling during start-up, with air in the generator enclosure first displaced by carbon dioxide before filling with hydrogen. This ensures that the highly flammable hydrogen does not mix with oxygen in the air.
The hydrogen pressure inside the casing is maintained slightly higher than atmospheric pressure to avoid outside air ingress, and up to about two atmospheres pressure to improve heat transfer capacity. The hydrogen must be sealed against outward leakage where the shaft emerges from the casing. Mechanical seals around the shaft are installed with a very small annular gap to avoid rubbing between the shaft and the seals on smaller turbines, with labyrinth type seals on larger machines.. Seal oil is used to prevent the hydrogen gas leakage to atmosphere.
The generator also uses water cooling. Since the generator coils are at a potential of about 22 kV, an insulating barrier such as Teflon is used to interconnect the water line and the generator high-voltage windings. Demineralized water of low conductivity is used.
Generator high-voltage system
The generator voltage for modern utility-connected generators ranges from in smaller units to in larger units. The generator high-voltage leads are normally large aluminium channels because of their high current as compared to the cables used in smaller machines. They are enclosed in well-grounded aluminium bus ducts and are supported on suitable insulators. The generator high-voltage leads are connected to step-up transformers for connecting to a high-voltage electrical substation (usually in the range of 115 kV to 765 kV) for further transmission by the local power grid.
The necessary protection and metering devices are included for the high-voltage leads. Thus, the steam turbine generator and the transformer form one unit. Smaller units may share a common generator step-up transformer with individual circuit breakers to connect the generators to a common bus.
Monitoring and alarm system
Most of the power station operational controls are automatic. However, at times, manual intervention may be required. Thus, the plant is provided with monitors and alarm systems that alert the plant operators when certain operating parameters are seriously deviating from their normal range.
Battery-supplied emergency lighting and communication
A central battery system consisting of lead–acid cell units is provided to supply emergency electric power, when needed, to essential items such as the power station's control systems, communication systems, generator hydrogen seal system, turbine lube oil pumps, and emergency lighting. This is essential for a safe, damage-free shutdown of the units in an emergency situation.
Circulating water system
To dissipate the thermal load of main turbine exhaust steam, condensate from gland steam condenser, and condensate from Low Pressure Heater by providing a continuous supply of cooling water to the main condenser thereby leading to condensation.
The consumption of cooling water by inland power stations is estimated to reduce power availability for the majority of thermal power stations by 2040–2069.
| Technology | Power generation | null |
2549891 | https://en.wikipedia.org/wiki/IceCube%20Neutrino%20Observatory | IceCube Neutrino Observatory | The IceCube Neutrino Observatory (or simply IceCube) is a neutrino observatory developed by the University of Wisconsin–Madison and constructed at the Amundsen–Scott South Pole Station in Antarctica. The project is a recognized CERN experiment (RE10).
Its thousands of sensors are located under the Antarctic ice, distributed over a cubic kilometer.
Similar to its predecessor, the Antarctic Muon And Neutrino Detector Array (AMANDA), IceCube consists of spherical optical sensors called Digital Optical Modules (DOMs), each with a photomultiplier tube (PMT) and a single-board data acquisition computer which sends digital data to the counting house on the surface above the array. IceCube was completed on 18 December 2010.
DOMs are deployed on strings of 60 modules each at depths between 1,450 and 2,450 meters into holes melted in the ice using a hot water drill. IceCube is designed to look for point sources of neutrinos in the teraelectronvolt (TeV) range to explore the highest-energy astrophysical processes.
Construction
IceCube is part of a series of projects developed and supervised by the University of Wisconsin–Madison. Collaboration and funding are provided by numerous other universities and research institutions worldwide. Construction of IceCube was only possible during the Antarctic austral summer from November to February, when permanent sunlight allows for 24-hour drilling. Construction began in 2005, when the first IceCube string was deployed and sufficient data was collected to verify that the optical sensors functioned correctly. In the 2005–2006 season, an additional eight strings were deployed, making IceCube the largest neutrino telescope in the world.
Construction was completed on 17 December 2010. The total cost of the project was $279 million.
As of 2024, plans for further upgrades to the array are in the federal approval process. If approved, the detectors for IceCube2 will each be eight times the size of those currently emplaced. The observatory will be able to detect more sources of particles, and discern their properties more finely at both lower and higher energy levels.
Sub-detectors
The IceCube Neutrino Observatory is composed of several sub-detectors which is also in addition to the main in-ice array.
AMANDA, the Antarctic Muon And Neutrino Detector Array, was the first part built, and it served as a proof-of-concept for IceCube. AMANDA was turned off in May 2009.
The IceTop array is a series of Cherenkov detectors on the surface of the glacier, with two detectors approximately above each IceCube string. IceTop is used as a cosmic ray shower detector, for cosmic ray composition studies and coincident event tests: if a muon is observed going through IceTop, it cannot be from a neutrino interacting in the ice.
The Deep Core Low-Energy Extension is a densely instrumented region of the IceCube array which extends the observable energies below 100 GeV. The Deep Core strings are deployed at the center (in the surface plane) of the larger array, deep in the clearest ice at the bottom of the array (between 1760 and 2450 m deep). There are no Deep Core DOMs between 1850 and 2107 m depth, as the ice is not as clear in those layers.
PINGU (Precision IceCube Next Generation Upgrade) is a proposed extension that will allow detection of low energy neutrinos (GeV energy scale), with uses including determining the neutrino mass hierarchy, precision measurement of atmospheric neutrino oscillation (both tau neutrino appearance and muon neutrino disappearance), and searching for WIMP annihilation in the Sun. A vision has been presented for a larger observatory, IceCube-Gen2.
Experimental mechanism
Neutrinos are electrically neutral leptons, and only interact very rarely with matter through the weak force. When they do react with the molecules of water in the ice via the charged current interaction, they create charged leptons (electrons, muons, or taus) corresponding to the flavor of the neutrino. These charged leptons can, if they are energetic enough, emit Cherenkov radiation. This happens when the charged particle travels through the ice faster than the speed of light in the ice, similar to the bow shock of a boat traveling faster than the waves it crosses. This light can then be detected by photomultiplier tubes within the digital optical modules making up IceCube.
The detector signatures of the three charged leptons are distinct, and as such it's possible to determine the neutrino flavor of charged current events. On the other hand if the neutrino scattered off the ice via the neutral current instead, the final state contains no information of the neutrino flavor since no charged lepton was created.
The signals from the PMTs are digitized and then sent to the surface of the glacier on a cable. These signals are collected in a surface counting house, and some of them are sent north via satellite for further analysis. Since 2014, hard drives rather than tape store the balance of the data which is sent north once a year via ship. Once the data reaches experimenters, they can reconstruct kinematical parameters of the incoming neutrino. High-energy neutrinos may cause a large signal in the detector, pointing back to their origin. Clusters of such neutrino directions indicate point sources of neutrinos.
Each of the above steps requires a certain minimum energy, and thus IceCube is sensitive mostly to high-energy neutrinos, in the range of 107 to about 1021 eV.
IceCube is more sensitive to muons than other charged leptons, because they are the most penetrating and thus have the longest tracks in the detector. Thus, of the neutrino flavors, IceCube is most sensitive to muon neutrinos. An electron resulting from an electron neutrino event typically scatters several times before losing enough energy to fall below the Cherenkov threshold; this means that electron neutrino events cannot typically be used to point back to sources, but they are more likely to be fully contained in the detector, and thus they can be useful for energy studies. These events are more spherical, or "cascade"-like, than "track"-like; muon neutrino events are more track-like.
Tau leptons can also create cascade events; but are short-lived and cannot travel very far before decaying, and are thus usually indistinguishable from electron cascades. A tau could be distinguished from an electron with a "double bang" event, where a cascade is seen both at the tau creation and decay. This is only possible with very high energy taus. Hypothetically, to resolve a tau track, the tau would need to travel at least from one DOM to an adjacent DOM (17 m) before decaying. As the average lifetime of a tau is , a tau traveling at near the speed of light would require 20 TeV of energy for every meter traveled. Realistically, an experimenter would need more space than just one DOM to the next to distinguish two cascades, so double bang searches are centered at PeV scale energies. Such searches are underway but have not so far isolated a double bang event from background events. Another way to detect lower energy tau neutrinos is through the "double pulse" signature, where a single DOM detect two distinct light arrival times corresponding to the neutrino interaction and tau decay vertices. One can also use machine learning (ML) techniques, such as Convolutional Neural Networks, to distinguish the tau neutrino signal. In 2024 the IceCube collaboration published its findings of seven astrophysical tau neutrino candidates using such a technique.
There is a large background of muons created not by neutrinos from astrophysical sources but by cosmic rays impacting the atmosphere above the detector. There are about 106 times more cosmic ray muons than neutrino-induced muons observed in IceCube. Most of these can be rejected using the fact that they are traveling downwards. Most of the remaining (up-going) events are from neutrinos, but most of these neutrinos are from cosmic rays hitting the far side of the Earth; some unknown fraction may come from astronomical sources, and these neutrinos are the key to IceCube point source searches. Estimates predict the detection of about 75 upgoing neutrinos per day in the fully constructed IceCube detector. The arrival directions of these astrophysical neutrinos are the points with which the IceCube telescope maps the sky. To distinguish these two types of neutrinos statistically, the direction and energy of the incoming neutrino is estimated from its collision by-products. Unexpected excesses in energy or excesses from a given spatial direction indicate an extraterrestrial source.
Experimental goals
Point sources of high energy neutrinos
A point source of neutrinos could help explain the mystery of the origin of the highest energy cosmic rays. These cosmic rays have energies high enough that they cannot be contained by galactic magnetic fields (their gyroradii are larger than the radius of the galaxy), so they are believed to come from extra-galactic sources. Astrophysical events which are cataclysmic enough to create such high energy particles would probably also create high energy neutrinos, which could travel to the Earth with very little deflection, because neutrinos interact so rarely. IceCube could observe these neutrinos: its observable energy range is about 100 GeV to several PeV. The more energetic an event is, the larger volume IceCube may detect it in; in this sense, IceCube is more similar to Cherenkov telescopes like the Pierre Auger Observatory (an array of Cherenkov detecting tanks) than it is to other neutrino experiments, such as Super-K (with inward-facing PMTs fixing the fiducial volume).
IceCube is more sensitive to point sources in the northern hemisphere than in the southern hemisphere. It can observe astrophysical neutrino signals from any direction, but neutrinos coming from the direction of the southern hemisphere are swamped by the cosmic-ray muon background. Thus, early IceCube point source searches focus on the northern hemisphere, and the extension to southern hemisphere point sources takes extra work.
Although IceCube is expected to detect very few neutrinos (relative to the number of photons detected by more traditional telescopes), it should have very high resolution with the ones that it does find. Over several years of operation, it could produce a flux map of the northern hemisphere similar to existing maps like that of the cosmic microwave background, or gamma ray telescopes, which use particle terminology more like IceCube. Likewise, KM3NeT could complete the map for the southern hemisphere.
IceCube scientists may have detected their first neutrinos on 29 January 2006.
Gamma-ray bursts coincident with neutrinos
When protons collide with one another or with photons, the result is usually pions. Charged pions decay into muons and muon neutrinos whereas neutral pions decay into gamma rays. Potentially, the neutrino flux and the gamma ray flux may coincide in certain sources such as gamma-ray bursts and supernova remnants, indicating the elusive nature of their origin. Data from IceCube is being used in conjunction with gamma-ray satellites like
Swift or Fermi for this goal. IceCube has not observed any neutrinos in coincidence with gamma ray bursts, but is able to use this search to constrain neutrino flux to values less than those predicted by the current models.
Indirect dark matter searches
Weakly interacting massive particle (WIMP) dark matter could be gravitationally captured by massive objects like the Sun and accumulate in the core of the Sun. With a high enough density of these particles, they would annihilate with each other at a significant rate. The decay products of this annihilation could decay into neutrinos, which could be observed by IceCube as an excess of neutrinos from the direction of the Sun. This technique of looking for the decay products of WIMP annihilation is called indirect, as opposed to direct searches which look for dark matter interacting within a contained, instrumented volume. Solar WIMP searches are more sensitive to spin-dependent WIMP models than many direct searches, because the Sun is made of lighter elements than direct search detectors (e.g. xenon or germanium). IceCube has set better limits with the 22 string detector (about of the full detector) than the AMANDA limits.
Neutrino oscillations
IceCube can observe neutrino oscillations from atmospheric cosmic ray showers, over a baseline across the Earth. It is most sensitive at ~25 GeV, the energy range for which the DeepCore sub-array has been optimized. DeepCore consists of 6 strings deployed in the 2009–2010 austral summer with a closer horizontal and vertical spacing. In 2014, DeepCore data was used to determine the mixing angle θ23 and mass splitting Δm223. This measurement has since been improved with more data and improved detector calibration and data processing.
As more data is collected and IceCube measurements are refined further, it may be possible to observe the characteristic modification of the oscillation pattern at ~15 GeV that determines the neutrino mass hierarchy. This mechanism for determining the mass hierarchy only works as the mixing angle θ13 is large.
Galactic supernovae
Despite the fact that individual neutrinos expected from supernovae have energies well below the IceCube energy cutoff, IceCube could detect a local supernova. It would appear as a detector-wide, brief, correlated rise in noise rates. The supernova would have to be relatively close (within our galaxy) to get enough neutrinos before the 1/r2 distance dependence took over. IceCube is a member of the Supernova Early Warning System (SNEWS).
Sterile neutrinos
A signature of sterile neutrinos would be a distortion of the energy spectrum of atmospheric neutrinos around 1 TeV, for which IceCube is uniquely positioned to search. This signature would arise from matter effects as atmospheric neutrinos interact with the matter of the Earth.
The described detection strategy, along with its South Pole position, could allow the detector to provide the first robust experimental evidence of extra dimensions predicted in string theory. Many extensions of the Standard Model of particle physics, including string theory, propose a sterile neutrino; in string theory this is made from a closed string. These could leak into extra dimensions before returning, making them appear to travel faster than the speed of light. An experiment to test this may be possible in the near future. Furthermore, if high energy neutrinos create microscopic black holes (as predicted by some aspects of string theory) it would create a shower of particles, resulting in an increase of "down" neutrinos while reducing "up" neutrinos.
In 2016, scientists at the IceCube detector did not find any evidence for the sterile neutrino.
Results
The IceCube collaboration has published flux limits for neutrinos from point sources, gamma-ray bursts, and neutralino annihilation in the Sun, with implications for WIMP–proton cross section.
A shadowing effect from the Moon has been observed. Cosmic ray protons are blocked by the Moon, creating a deficit of cosmic ray shower muons in the direction of the Moon. A small (under 1%) but robust anisotropy has been observed in cosmic ray muons.
In November 2013 it was announced that IceCube had detected 28 neutrinos that likely originated outside the Solar System and among those a pair of high energy neutrinos in the peta-electron volt range, making them the highest energy neutrinos discovered to date.
The pair were nicknamed "Bert" and "Ernie", after characters from the Sesame Street TV show.
Later in 2013 the number of detection increased to 37 candidates including a new high energy neutrino at 2000-TeV given the name of "Big Bird".
IceCube measured 10–100 GeV atmospheric muon neutrino disappearance in 2014, using three years of data taken May 2011 to April 2014, including DeepCore, determining neutrino oscillation parameters ∆m232 = × 10−3 eV2 and sin2(θ23) = (normal mass hierarchy), comparable to other results. The measurement was improved using more data in 2017, and in 2019 atmospheric tau neutrino appearance was measured. The latest measurement with improved detector calibration and data processing from 2023 has resulted in more precise values of the oscillation parameters, determining ∆m232 = (2.41 ± 0.07) × 10−3 eV2 and sin2(θ23) = 0.51 ± 0.05 (normal mass hierarchy).
In July 2018, the IceCube Neutrino Observatory announced that they had traced an extremely-high-energy neutrino that hit their detector in September 2017 back to its point of origin in the blazar TXS 0506 +056 located 5.7 billion light-years away in the direction of the constellation Orion, the results had a statistical significance of 3-3.5σ. This was the first time that a neutrino detector had been used to locate an object in space, and indicated that a source of cosmic rays had been identified.
In 2020, evidence of the Glashow resonance at 2.3σ (formation of the W boson in antineutrino-electron collisions) was announced.
In February 2021, the tidal disruption event (TDE) AT2019dsg was reported as candidate for a neutrino source and the TDE AT2019fdr as a second candidate in June 2022.
In November 2022, IceCube announced strong evidence of a neutrino source emitted by the active galactic nucleus of Messier 77. It is the second detection by IceCube after TXS 0506+056, and only the fourth known source including SN1987A and solar neutrinos. OKS 1424+240 and GB9 are other possible candidates.
In June 2023 IceCube identified as a galactic map the neutrino diffuse emission from the Galactic plane at the 4.5σ level of significance.
| Technology | Ground-based observatories | null |
2552693 | https://en.wikipedia.org/wiki/Xerini | Xerini | Xerini is a tribe of ground squirrels occurring in Africa and Asia. With the tribes Marmotini (Holarctic ground squirrels) and Protoxerini (African tree squirrels), they form the subfamily Xerinae. There are five living genera—Xerus, the unstriped ground squirrel; Euxerus, the striped ground squirrel; Geosciurus, the Cape and mountain ground squirrels; Atlantoxerus, containing the living Barbary ground squirrel of North Africa and some extinct species; and Spermophilopsis, containing the long-clawed ground squirrel of Central Asia.
The squirrels live in open woodlands, grasslands, or rocky country. They are diurnal and terrestrial, living in burrows. Their diet is roots, seeds, fruits, pods, grains, insects, small vertebrates and bird eggs. They live in colonies similar to North American prairie dogs, and have similar behavior. If kept as pets, they run free as house cats do, otherwise farmers consider them pests. Breeding in South African ground squirrels is asynchronous and there is no specific period of breeding although very few litters are seen above the ground in the months of July to October. Gestation period lasts for 48 days and the young ones are weaned after 52 days. There are one to three babies per litter. A female becomes sexually mature when she is 10 months old and a male matures at the age of 8 months. Although a female has the capability to breed throughout the year, less than 10% reproduce more than one litter in a year.
The South African ground squirrel does not hibernate. These squirrels are very social and live in groups with about 1 to 3 females and 2 to 3 males. Sometimes, the number of males can exceed up to 9 with some sub-adult females. An interesting thing about this breed is that mature males like to form their own groups and there can be 19 to 20 individuals in one group. These groups are called bands.
Their home range is shared with other groups which is actually quite a surprising thing, since in case of most of the mammals there is generally an extremely tough competition to get the female. The group-living instinct of this species gives it a survival opportunity from predators overweighing the disadvantage of breeding opportunities which come from competition.
Species
Extant species
Genus Atlantoxerus
Barbary ground squirrel, Atlantoxerus getulus
Genus Euxerus
Striped ground squirrel, Euxerus erythropus
Genus Geosciurus
Cape ground squirrel, Geosciurus inauris
Mountain ground squirrel, Geosciurus princeps
Genus Xerus
Unstriped ground squirrel, Xerus rutilus
Genus Spermophilopsis
Long-clawed ground squirrel, Spermophilopsis leptodactylus
Fossil species
Genus Heteroxerus
†Heteroxerus grivensis
Genus Aragoxerus
†Aragoxerus ignis
| Biology and health sciences | Rodents | Animals |
2552745 | https://en.wikipedia.org/wiki/Sparassodonta | Sparassodonta | Sparassodonta (from Greek [], to tear, rend; and , gen. [, ], tooth) is an extinct order of carnivorous metatherian mammals native to South America, related to modern marsupials. They were once considered to be true marsupials, but are now thought to be a separate side branch that split before the last common ancestor of all modern marsupials.
A number of these mammalian predators closely resemble placental predators that evolved separately on other continents, and are cited frequently as examples of convergent evolution. They were first described by Florentino Ameghino, from fossils found in the Santa Cruz beds of Patagonia. Sparassodonts were present throughout South America's long period of "splendid isolation" during the Cenozoic; during this time, they shared the niches for large warm-blooded predators with the flightless terror birds. Previously, it was thought that these mammals died out in the face of competition from "more competitive" placental carnivorans during the Pliocene Great American Interchange, but more recent research has showed that sparassodonts died out long before eutherian carnivores arrived in South America (aside from procyonids, which sparassodonts probably did not directly compete with).
Sparassodonts have been referred to as borhyaenoids by some authors, but currently the term Borhyaenoidea refers to a restricted subgroup of sparassodonts comprising borhyaenids and their close relatives.
Anatomy
Almost all sparassodonts have an exceptionally shortened snout—most especially thylacosmylids. Hathliacynids usually have a longer snout than the other groups. The nasal bones extend past the eye sockets, often reaching the lacrimal bone. Except for thylacosmylids beyond Patagosmilus, sparassodonts feature an open eye socket, with more marginalized (though nonetheless prominent) postorbital processes which would otherwise form the postorbital bar connecting the forehead to the cheek, thus framing the eye. They exhibit marked postorbital constriction. The orbital process (between the cheek and the eye socket) is usually diminished, though the zygomatic arch (the cheekbone) is strong. They feature a prominent sagittal crest along the midline of the flattened skull, the crest strength is quite variable among borhyaenids. They have an expanded occipital bone with a well defined nuchal crest.
Sparassodonts spanned a wide range of body sizes, from 2.2-pound (1 kg) weasel or civet-like forms to Thylacosmilus, which was the size of a leopard. Along with the Australian thylacoleonids, sparassodonts include some of the largest metatherian carnivores.
Sparassodonts have highly reduced epipubic bones (pelvic bones which support the pouch), to the point that early analysis could not even find evidence for them. This is a characteristic shared with the Australian thylacine, and historically argued as a synapomorphy, though nowadays it is considered to have developed independently for poorly understood reasons. As with thylacines, it is very likely that they possessed long cartilaginous elements instead.
Teeth
The dental formula of sparassodonts varies considerably. In borhyaenids, it is , with three upper and lower incisors, one upper and lower canine, three upper and lower premolars, and four upper and lower molars in each half of either jaw. Proborhyaenids usually only have two lower incisors instead of three, except for Callistoe. Thylacosmylids have at least two upper and only two lower incisors (the uppers grew into elongated sabers), and two upper and lower premolars. Some specimens of Borhyaena and Arctodictis are also missing the last upper molar, showing that the presence of this tooth was variable in these species.
Sparassodonta is characterized by dental synapomorphies that distinguish the group from other closely related mammals. Unequivocal traits uniting the earliest Sparassodonts include:
a snout that forms a pronounced bulge around the canine teeth when viewed from above
a ridge on the upper molar (preparacrista) oriented anterobuccally (towards the cheek) with respect to the long axis of the tooth.
a pronounced keel near the base of the front of the paraconid
ridges on lower molars (postprotocristid-metacristid) parallel or oblique with respect to lower jaw axis.
a very tall protoconid (>90% tooth length, secondarily lost in Hondadelphys and Stylocynus) that bulges to the side and is wider at its midpoint than its base
talonid (crushing end) of lower molar narrow in relation to trigonid (shearing end).
In borhyaenids, only the third premolar was ever replaced in the animal's lifetime, similar to other metatherians. In thylacosmilids, only the lower third premolar was replaced.
The cusps of the sparassodont molar correlate to a cutting function rather than a crushing one. In the upper molars, the paracone (on the lingual, or tongueward, side) is reduced and fused to the metacone (distal, towards the back of the mouth), inflating the postmetacrista (the lingual border of the metacone); and they almost always lack the stylar shelf (on the buccal, or cheekward, side) and associated stylar cusps. In the lower molars, the trigonids (the buccal shearing side) have an inflated paracristid and marginalized or absent metaconid; and the talonid (the distal, or backendwards, crushing side) is either reduced or gone.
Taxonomy
Classification
Sparassodonts can be divided into six major groups; basal sparassodonts (?earliest Paleocene-late Miocene), species that cannot be easily assigned to any of the other sparassodont groups and whose teeth often exhibit adaptations for omnivory; hathliacynids (late Oligocene-early Pliocene/late Pliocene), which range from a marten to a thylacine in size, and have long, fox-like muzzles and teeth strongly suited for carnivory; basal borhyaenoids (middle Eocene-late Miocene), borhyaenoids which are unable to be easily classified into the families Borhyaenidae, Thylacosmilidae, or Proborhyaenidae and range in form and size; borhyaenids (early-late Miocene), the sparassodont group most specialized for running, but not as much as living carnivorans or even thylacines; proborhyaenids (middle Eocene-late Oligocene), robust, wolverine-like forms with ever-growing upper and lower canines; and thylacosmilids (early Miocene-late Pliocene), another terrestrially specialized group with ever-growing saber-like upper canines.
The taxonomic classification below follows the latest review of the group, that of Prevosti and Forasiepi (2018), with additions from more recent studies. Although Mayulestes was originally described as a sparassodont, later phylogenetic analyses have shown that it most likely does not belong to this group; however more recent studies show it to be closely related to sparassodonts. Similarly, while basal borhyaenoids such as Lycopsis and Prothylacynus were once thought to belong to a distinct family (Prothylacynidae), phylogenetic analyses have found that these animals do not represent a monophyletic group. The exact age of most Eocene species of sparassodonts is uncertain, given the lack of precise stratigraphic information associated with most specimens and the recent division of the Casamayoran SALMA into the Vacan and Barrancan SALMAs.
Order Sparassodonta
Genus Allqokirus
Allqokirus australis (earliest Paleocene, Tiupampan SALMA)
Genus Mayulestes
Mayulestes ferox (earliest Paleocene, Tiupampan SALMA)
Genus Argyrolestes
Genus Nemolestes
Nemolestes spalacotherinus (late middle Eocene, Barrancan SALMA)
Genus Patene
Patene campbelli (late ?Eocene, Santa Rosa local fauna)
Patene coloradensis (middle Eocene)
Patene coluapiensis (middle Eocene, Barrancan SALMA)
Patene simpsoni (late Paleocene/early Eocene, Itaboraian SALMA)
Genus Procladosictis
Procladosictis anomala (late Eocene, Mustersan SALMA)
Family Hondadelphidae
Genus Hondadelphys
Hondadelphys fieldsi (late middle Miocene, Laventan SALMA)
Genus Stylocynus
Stylocynus paranensis (latest Miocene, Huayquerian SALMA)
Family Hathliacynidae
Genus Acyon
Acyon ?herrerae (early Miocene, Colhuehuapian SALMA)
Acyon myctoderos (late middle Miocene, Laventan SALMA)
Acyon tricuspidatus (late early Miocene, Santacrucian SALMA)
Genus Australogale
Australogale leptognathus (late middle Miocene, Laventan SALMA)
Genus Borhyaenidium
Borhyaenidium altiplanicus (latest Miocene, Huayquerian SALMA)
Borhyaenidium riggsi (unknown, either early or late Pliocene, Montehermosan or Chapadmalalan SALMA)
Borhyaenidium musteloides (latest Miocene, Huayquerian SALMA)
Genus Chasicostylus
Chasicostylus castroi (early late Miocene, Chasicoan SALMA)
Genus Cladosictis
Cladosictis centralis (early Miocene, Colhuehuapian SALMA)
Cladosictis patagonica (late early to earliest middle Miocene, Santacrucian-Friasian SALMA)
Genus Notictis
Notictis ortizi (latest Miocene, Huayquerian SALMA)
Genus Notocynus
Notocynus hermosicus (early Pliocene, Montehermosan SALMA)
Genus Notogale
Notogale mitis (late Oligocene, Deseadan SALMA)
Genus Pseudonotictis
Pseudonotictis chubutensis (early middle Miocene, Colloncuran SALMA)
Pseudonotictis pusillus (late early Miocene, Santacrucian SALMA)
Genus Perathereutes
Perathereutes pungens (late early Miocene, Santacrucian SALMA)
Genus Sallacyon
Sallacyon hoffstetteri (late Oligocene, Deseadan SALMA)
Genus Sipalocyon
Sipalocyon externus (early Miocene, Colhuehuapian SALMA)
Sipalocyon gracilis (late early Miocene, Santacrucian SALMA)
Sipalocyon "obusta" (late early Miocene, Santacrucian SALMA)
Superfamily Borhyaenoidea
Genus Angelocabrerus
Angelocabrerus daptes (middle Eocene)
Genus Chlorocyon
Chlorocyon phantasma (late Eocene, Mustersan SALMA)
Genus Dukecynus
Dukecynus magnus (late middle Miocene, Laventan SALMA)
Genus Fredszalaya
Fredszalaya hunteri (late Oligocene, Deseadan SALMA)
Genus Lycopsis
Lycopsis longirostrus (late middle Miocene, Laventan SALMA)
Lycopsis padillai (middle Miocene, Colloncuran SALMA)
Lycopsis torresi (early Miocene, Colhuehuapian SALMA)
Lycopsis viverensis (early late Miocene, Chasicoan SALMA))
Genus Pharsophorus
Pharsophorus lacerans (middle to late Oligocene, La Cantera Local Fauna to Deseadan SALMA)
Pharsophorus tenax (late Oligocene, Deseadan SALMA)
Genus Plesiofelis
Plesiofelis schlosseri (late Eocene, Mustersan SALMA)
Genus Prothylacynus
Prothylacynus patagonicus (late early to earliest middle Miocene, Santacrucian-Friasian SALMA)
Genus Pseudolycopsis
Pseudolycopsis cabrerai (early late Miocene, Chasicoan SALMA
Genus Pseudothylacynus
Pseudothylacynus rectus (early Miocene, Colhuehuapian SALMA)
Family Borhyaenidae
Genus Acrocyon
Acrocyon riggsi (early Miocene, Colhuehuapian SALMA)
Acrocyon sectorius (late early Miocene, Santacrucian SALMA)
Genus Arctodictis
Arctodictis munizi (late early Miocene, Santacrucian SALMA)
Arctodictis sinclairi (early Miocene, Colhuehuapian SALMA)
Genus Australohyaena
Australohyaena antiquua (late Oligocene, Deseadan SALMA)
Genus Borhyaena
Borhyaena macrodonta (early Miocene, Colhuehuapian SALMA)
Borhyaena tuberata (late early Miocene, Santacrucian SALMA)
?Genus Eutemnodus
Family Proborhyaenidae
Genus Arminiheringia
Genus Callistoe
Callistoe vincei (middle Eocene, Vacan-Barrancan SALMA)
Genus Paraborhyaena
Paraborhyaena boliviana (late Oligocene, Deseadan SALMA)
Genus Proborhyaena
Proborhyaena gigantea (late Oligocene, Deseadan SALMA)
Subfamily Thylacosmilinae
Genus Eomakhaira
Eomakhaira molossus (early Oligocene, ?Tinguirirican SALMA)
Genus Anachlysictis
Anachlysictis gracilis (late middle Miocene, Laventan SALMA)
Genus Patagosmilus
Patagosmilus goini (early middle Miocene, Colloncuran SALMA)
Genus Thylacosmilus
Thylacosmilus atrox (latest Miocene to late Pliocene, Huayquerian-Chapadmalalan SALMA)
Several other metatherian taxa have been suggested to be sparassodonts or closely related to sparassodonts. The australian Murgon taxa Archaeonothos has been noted as being similar to sparassodonts, but currently its relationships are not fully concluded. Carneiro (2018) recovered the genus Varalphadon from the Late Cretaceous of North America as a basal member of Sparassodonta. However, this interpretation of Varalphadon as a sparassodont has not been supported by later phylogenetic analyses, and most of the purported synapomorphies between Varalphadon and sparassodonts are not actually present in Varalphadon or have been suggested to be due to convergent evolution. Sparassodonts have sometimes been considered closely related to the "Gurlin Tsav skull" an unnamed metatherian known from a partial skull found in the Late Cretaceous Nemegt Formation of Mongolia.
The following cladogram of sparassodont interrelationships is after Engelman et al., 2020. Not all studies agree on the sister group relationship between Thylacosmilidae and Borhyaenidae recovered here, with other studies finding thylacosmilids to be within Proborhyaenidae. The relationships among hathliacynids are also relatively unstable.
Within Metatheria, a 2016 phylogenetic analysis group found that borhyaenids form a clade with the Asian "Gurlin Tsav skull" as well as other South American taxa. The same phylogeny found that marsupials group among various North American Cretaceous species. The phylogenetic tree is reproduced below.
Evolution
The early history of the Sparassodonta is poorly known, as most Paleocene and Eocene members of this group are only known from isolated teeth and fragmentary jaws. However, one species, the middle Eocene Callistoe vincei, is known from a nearly complete, articulated skeleton. As Callistoe belongs to one of the most specialized groups of sparassodonts, this indicates that the other major groups (e.g. borhyaenids, hathliacynids, etc.) must have also arisen by this time. Originally, the early Paleocene metatherian Mayulestes was considered to be the earliest known member of the Sparassodonta, but phylogenetic analyses suggest that this species represents an independent radiation of carnivorous metatherians more closely related to Pucadelphys; however, recent studies show that these taxa were closely related to borhyaenids. As of this writing, the earliest known true sparassodonts are either Allqokirus australis, a species from the same site as Mayulestes that may turn out to not be a sparassodont, and an isolated astragalus from the earliest Paleocene site of Punta Peligro, Argentina.
Sparassodonts have been suggested to be related to a variety of other groups of metatherians. Florentino Ameghino, who first described fossils of the group, thought that sparassodonts were closely related to creodonts and were a transitional group between metatherians and carnivorous placentals (including modern carnivorans). Contemporary authors in the late 19th and early 20th century rejected this hypothesis and considered sparassodonts to be closely related to Australian thylacines and dasyurids. The most popular hypothesis for much of the 20th century was that sparassodonts were closely related to opossums. In 1990, Marshall et al. (1990) considered the Cretaceous stagodontids to be members of Sparassodonta, but this was criticized by later authors. Marshall and Kielan-Jaworowska (1992) considered sparassodonts to be closely related to deltatheroidans, but this was also criticized. Most of these hypotheses were based on similar adaptations for carnivorous diets in sparassodonts, opossums, dasyuromorphians, stagodonts, and deltatheroidans, which are highly prone to convergent evolution within mammals. Szalay (1994) considered sparassodonts to be closely related to paucituberculatans based on features of the ankle. In recent years there has been a growing consensus that sparassodonts are positioned just outside of crown-group Marsupialia, in a broader clade (Pucadelphyida) including pucadelphyids as well as sparassodonts.
Sparassodonts are presently regarded as an endemic South American group, and have not even been found in nearby continents like Antarctica (where other groups native to South America such as litopterns, astrapotheres, microbiotheres, and polydolopids) are present.
Paleobiology
Diet
Sparassodonts were carnivorous, and with the exception of some basal members of all members of this group were hypercarnivorous (having diets composed of more than 70% meat). Only Hondadelphys and Stylocynus appear to have exhibited adaptations for omnivory, and even then Stylocynus may have had a more mesocarnivorous diet similar to canids than an omnivorous one. Medium-to-large caviomorph rodents and rodent-like mammals (e.g., small notoungulates) appear to have been common prey items of sparassodonts. The subadult holotype of Lycopsis longirostrus preserves remains of the dinomyid Scleromys colombianus as fossilized gut contents. Sparassodont coprolites from the Santa Cruz Formation preserve the bones of chinchillid and octodontoid rodents inside them. Bite marks from medium-sized sparassodonts have been found on the small notoungulate Paedotherium. Stable isotope data from the early late Miocene Lycopsis viverensis and Thylacosmilus atrox suggests that these species fed on C3 grazers in open habitats, likely notoungulates.
Bite marks likely pertaining to hathliacynid sparassodonts have been found on the remains of penguins and flightless marine ducks in ancient seabird nesting colonies, suggesting that sparassodonts raided seabird colonies for eggs, carrion, and other prey like many predatory mammals do today.
Borhyaenid and proborhyaenid sparassodonts have been interpreted as being capable of crushing bones similar to modern hyenas, wolverines, or the Tasmanian devil (Sarcophilus harrisii) based on their deep jaws, bulbous premolars with deep roots and pronounced wear at their tips, extensive fused or interlocking mandibular symphyses, large masseteric fossae, microfractures in their tooth enamel, and high estimated bite forces. Australohyaena antiquua shows particularly pronounced adaptations for bone-cracking, with a very deep jaw and strongly arched nasals similar to what is seen in modern hyaenids.
Based on studies of the postcranial skeleton, it appears as though most sparassodonts were scansorial (adapted for climbing), although terrestrial adaptations evolved in Lycopsis longirostrus, borhyaenids, proborhyaenids, and thylacosmilids. Most sparassodonts were plantigrade, Borhyaena has been suggested to have been digitigrade but this has been questioned. The one exception was Thylacosmilus, which has been interpreted as having a digitigrade forefoot and a semiplantigrade hindfoot, this has been supported by fossil tracks.
One unusual aspect of sparassodont paleoecology is that at most fossil localities their remains are nearly ten times rarer than would be expected based on comparisons with carnivorous mammals at fossil sites in other parts of the world. The exact reasons for this are not clear, though this appears to be a broader pattern applicable to other groups of Cenozoic South American terrestrial carnivores (i.e., terror birds).
Sociality
Little is known of the behavior and biology of sparassodonts outside of general locomotor and dietary habits. Argot (2004) proposed that Thylacosmilus atrox may have exhibited protracted parental care after weaning of the offspring, given that saber teeth in general have been suggested to require long juvenile periods for the young to gain the skill necessary to use them effectively. However, this has not been tested further. Sparassodonts have relatively large and complex brains for metatherians, comparable to those of some Australian marsupials like Australian possums, though the body masses used to produce these estimates of relative brain size are low compared later studies suggesting these values could be overestimated.
Wounds have been documented on the face of specimens of Borhyaena tuberata and Sipalocyon gracilis, potentially suggesting aggressive habits similar to the modern Tasmanian devil (Sarcophilus harrisii).
Senses
Sparassodonts appear to have had very little binocular vision, with borhyaenids having the greatest degree of depth perception (but still lower than modern carnivorans) and the eyes of Thylacosmilus facing almost completely to the sides. However, later studies have found that Thylacosmilus likely held its head in a downward-facing position, which would have allowed for more binocular vision than previously thought.
Pathology
Several specimens of hathliacynids (Sipalocyon and Cladosictis) show a pathological disorder characterized by the presence of growths on the surface of the mandible, which in the most extreme cases can result in the loss of several teeth due to bony pathological growths. The exact cause of this condition (i.e., infection, virus, parasite) and why it seems to only occur in small sparassodonts is unknown, though this condition has also been documented in microbiotherians.
Extinction
After the middle Miocene, sparassodonts began to slowly decline in diversity. Basal borhyaenoids are last known from the early late Miocene (Pseudolycopsis cabrerai and Lycopsis viverensis), and after this time were at least partially replaced by large-bodied basal sparassodonts such as Stylocynus. Some have proposed that this shift in dominance was because of the more omnivorous habits of basal sparassodonts, which may have been more adapted to the more seasonal South American climates of the late Neogene. Borhyaenids are last known from the latest Miocene, though only fragmentary remains of this group are known from this period. Later remains assigned to this group have since been reidentified as thylacosmilids or procyonids. By the Pliocene, only two families of sparassodonts remained in South America, the Hathliacynidae and the Thylacosmilidae. Pliocene hathliacynid remains are rare, and it is possible that these animals may have competed with the large carnivorous didelphids such as Lutreolina that appeared around this time. Hathliacynids are last definitively known from the early Pliocene, though their remains are rare.
The thylacosmilids, on the other hand, were more successful and abundant, being some of the only large mammalian carnivores in South America during the Pliocene, before dying out during a faunal turnover in the middle of the epoch (the youngest specimens of thylacosmilids are ~3.3 Ma). It is still not certain why Sparassodonta declined in diversity and became extinct during the late Cenozoic, but it appears as though competition from eutherian carnivorans was not a factor, as the placental analogues of sparassodonts (dogs, weasels, and saber-toothed cats) did not enter South America until the middle Pleistocene, several million years after their sparassodont counterparts became extinct. Sparassodonts did coexist with Cyonasua-group procyonids during the late Miocene and Pliocene, but Cyonasua-group procyonids appear to have been primarily omnivorous and filled ecological niches that sparassodonts never occupied, which may be one reason that these animals were able to colonize South America despite the diverse predator guild in the late Miocene. The overall decline in sparassodont diversity from the Late Miocene to the end of the Pliocene may be linked to the climatic cooling that characterised the Late Neogene and the onset of the Quaternary glaciation. Additionally, the increased aridity of South America caused by the uplift of the Andes was likely responsible as well.
| Biology and health sciences | Marsupials | Animals |
2552781 | https://en.wikipedia.org/wiki/Phorusrhacos | Phorusrhacos | Phorusrhacos ( ) is an extinct genus of giant flightless terror birds that inhabited South America during the Miocene epoch. Phorusrhacos was one of the dominant land predators in South America at the time it existed. It is thought to have lived in woodlands and grasslands.
Discovery and naming
Remains are known from several localities in the Santa Cruz Formation and Monte León Formation in Santa Cruz Province, of Argentina. Among the bones found in the strata of the Santa Cruz Formations (now considered as mainly of mid-Miocene date) was the piece of a mandible which Florentino Ameghino discovered in early 1887 and the same year at first described as that of an edentate mammal which he named Phorusrhacos longissimus.
The generic name is derived from Greek -φόρος, (-phoros), an element meaning "bearer" in word combinations, and ῥάκος, (rhakos), "rag" or "wrinkle", probably in reference to the wrinkled jaw surface. When the original derivation was no longer understood, other translations were given, such as the literal translation of "Rag-Thief", and "branch-holder" from the mistaken assumption the name had been intended to be derived from a Greek rhakis, "branch". The specific name means "very long" in Latin, again in reference to the lower jaws. The holotype is the mandible, specimen MLP-118 (Museo de La Plata). In 1889 Ameghino emended the name to a more grammatically correct Phororhacos but the earlier name has priority. In 1891, it was by him recognized to be a bird.
Description
Phorusrhacos had a skull nearly long, stood nearly tall, and probably weighed nearly , as much as a male ostrich. It had very strong legs, capable of running at high speed, stubby, flightless wings, a long neck, and a proportionately large head. This ended in a huge, hooked beak that could tear through flesh easily, or stab into prey. The lower jaw was smaller than the upper jaw. There were three toes on each of the feet, all of which were armed with sharp claws.
Classification
Phorusrhacos was part of the group called the Phorusrhacidae, which is an extinct group of flightless, cursorial carnivorous birds that occupied one of the dominant, large land-predator niches in South America from the lower Eocene to the Pleistocene. They dispersed into North America during the Great American Biotic Interchange (~3 Ma). Some remains from Africa and Europe and the Paleocene of Brazil have been referred to this clade or identified as phylogenetically related to the extant South American seriemas, but these assignments remain controversial.
The following cladogram follows the analysis of Degrange and colleagues, 2015:
| Biology and health sciences | Prehistoric birds | Animals |
1836606 | https://en.wikipedia.org/wiki/Allylic%20rearrangement | Allylic rearrangement | An allylic rearrangement or allylic shift is an organic chemical reaction in which reaction at a center vicinal to a double bond causes the double bond to shift to an adjacent pair of atoms:
It is encountered in both nucleophilic and electrophilic substitution, although it is usually suppressed relative to non-allylic substitution. For example, reaction of 1-chloro-2-butene with sodium hydroxide gives 2-buten-1-ol and 3-buten-2-ol:
In the similar substitution of 1-chloro-3-methyl-2-butene, the secondary 2-methyl-3-buten-2-ol is produced in a yield of 85%, while that for the primary 3-methyl-2-buten-1-ol is 15%.
Allylic shifts occur because the transition state is an allyl intermediate. In other respects they are similar to classical nucleophilic substitution, and admit both bimolecular and monomolecular mechanisms (respectively the SN2' and SN1'/SNi' substitutions).
Scope
Allylic shifts become the dominant reaction pathway when there is substantial resistance to a normal (non-allylic) substitution. For nucleophilic substitution, such resistance is known when there is substantial steric hindrance at or around the leaving group, or if there is a geminal substituent destabilizing an accumulation of positive charge. The effects of substitution at the vinyl group are less clear.
Although rarer still than SN', allylic shifts can occur vinylogously, as a "butadienylic shift":
SN2' reduction
In SN2' reduction, a hydride allylically displaces a good leaving group in a formal organic reduction, similar to the Whiting diene synthesis. One example occurred in taxol total synthesis (ring C):
The hydride is lithium aluminium hydride and the leaving group a phosphonium salt; the allylic shift causes the exocyclic double bond in the product. Only when the cyclohexane ring is properly substituted will the proton add trans to the adjacent methyl group.
Electrophilic allyl shifts
Allyl shifts can also take place with electrophiles. In the example below the carbonyl group in benzaldehyde is activated by diboronic acid prior to reaction with the allyl alcohol (see: Prins reaction):
The active catalyst system in this reaction is a combination of a palladium pincer compound and p-toluenesulfonic acid, the reaction product is obtained as a single regioisomer and stereoisomer.
Examples
Repeated allylic shifts can "flip-flop" a double-bond between two possible locations:
An SN2' reaction should explain the outcome of the reaction of an aziridine carrying a methylene bromide group with methyllithium:
In this reaction one equivalent of acetylene is lost.
Named reactions
Ferrier rearrangement
Meyer–Schuster rearrangement
| Physical sciences | Organic reactions | Chemistry |
1837735 | https://en.wikipedia.org/wiki/Wittig%20reaction | Wittig reaction | The Wittig reaction or Wittig olefination is a chemical reaction of an aldehyde or ketone with a triphenyl phosphonium ylide called a Wittig reagent. Wittig reactions are most commonly used to convert aldehydes and ketones to alkenes. Most often, the Wittig reaction is used to introduce a methylene group using methylenetriphenylphosphorane (Ph3P=CH2). Using this reagent, even a sterically hindered ketone such as camphor can be converted to its methylene derivative.
Reaction mechanism
Mechanistic studies have focused on unstabilized ylides, because the intermediates can be followed by NMR spectroscopy. The existence and interconversion of the betaine (3a and 3b) is subject of ongoing research. For lithium-free Wittig reactions, studies support a concerted formation of the oxaphosphetane without intervention of a betaine. In particular, phosphonium ylides 1 react with carbonyl compounds 2 via a [2+2] cycloaddition that is sometimes described as having [π2s+π2a] topology to directly form the oxaphosphetanes 4a and 4b. Under lithium-free conditions, the stereochemistry of the product 5 is due to the kinetically controlled addition of the ylide 1 to the carbonyl 2. When lithium is present, there may be equilibration of the intermediates, possibly via betaine species 3a and 3b. Bruce E. Maryanoff and A. B. Reitz identified the issue about equilibration of Wittig intermediates and termed the process "stereochemical drift". For many years, the stereochemistry of the Wittig reaction, in terms of carbon-carbon bond formation, had been assumed to correspond directly with the Z/E stereochemistry of the alkene products. However, certain reactants do not follow this simple pattern. Lithium salts can also exert a profound effect on the stereochemical outcome.
Mechanisms differ for aliphatic and aromatic aldehydes and for aromatic and aliphatic phosphonium ylides. Evidence suggests that the Wittig reaction of unbranched aldehydes under lithium-salt-free conditions do not equilibrate and are therefore under kinetic reaction control. E. Vedejs has put forth a theory to explain the stereoselectivity of stabilized and unstabilized Wittig reactions.
Strong evidence indicated that under Li-free conditions, Wittig reactions involving unstabilized (R1= alkyl, H), semistabilized (R1 = aryl), and stabilized (R1 = EWG) Wittig reagents all proceed via a [2+2]/retro-[2+2] mechanism under kinetic control, with oxaphosphetane as the one and only intermediate.
Scope and limitations
Functional group tolerance
The Wittig reagents generally tolerate carbonyl compounds containing several kinds of functional groups such as OH, OR, nitroarenes, epoxides, and sometimes esters and amides. Even ketone, aldehyde, and nitrile groups can be present if conjugated with the ylide — these are the stabilised ylides mentioned above. Bis-ylides (containing two P=C bonds) have also been made and used successfully. There can be a problem with sterically hindered ketones, where the reaction may be slow and give poor yields, particularly with stabilized ylides, and in such cases the Horner–Wadsworth–Emmons (HWE) reaction (using phosphonate esters) is preferred. Another reported limitation is the often labile nature of aldehydes, which can oxidize, polymerize or decompose. In a so-called tandem oxidation-Wittig process the aldehyde is formed in situ by oxidation of the corresponding alcohol.
Stereochemistry
For the reaction with aldehydes, the double bond geometry is readily predicted based on the nature of the ylide. With unstabilised ylides (R3 = alkyl) this results in (Z)-alkene product with moderate to high selectivity. If the reaction is performed in dimethylformamide in the presence of lithium iodide or sodium iodide, the product is almost exclusively the Z-isomer. With stabilized ylides (R3 = ester or ketone), the (E)-alkene is formed with high selectivity. The (E)/(Z) selectivity is often poor with semistabilized ylides (R3 = aryl).
To obtain the (E)-alkene for unstabilized ylides, the Schlosser modification of the Wittig reaction can be used. Alternatively, the Julia olefination and its variants also provide the (E)-alkene selectively. Ordinarily, the Horner–Wadsworth–Emmons reaction provides the (E)-enoate (α,β-unsaturated ester), just as the Wittig reaction does. To obtain the (Z)-enolate, the Still-Gennari modification of the Horner-Wadsworth-Emmons reaction can be used.
Schlosser modification
The main limitation of the traditional Wittig reaction is that the reaction proceeds mainly via the erythro betaine intermediate, which leads to the Z-alkene. The erythro betaine can be converted to the threo betaine using phenyllithium at low temperature. This modification affords the E-alkene.
Allylic alcohols can be prepared by reaction of the betaine ylide with a second aldehyde. For example:
Example
An example of its use is in the synthesis of leukotriene A methyl ester. The first step uses a stabilised ylide, where the carbonyl group is conjugated with the ylide preventing self condensation, although unexpectedly this gives mainly the cis product. The second Wittig reaction uses a non-stabilised Wittig reagent, and as expected this gives mainly the cis product.
History
The Wittig reaction was reported in 1954 by Georg Wittig and his coworker Ulrich Schöllkopf. In part for this contribution, Wittig was awarded the Nobel Prize in Chemistry in 1979.
| Physical sciences | Organic reactions | Chemistry |
1838808 | https://en.wikipedia.org/wiki/Needlefish | Needlefish | Needlefish (family Belonidae) or long toms are piscivorous fishes primarily associated with very shallow marine habitats or the surface of the open sea. Some genera include species found in marine, brackish, and freshwater environments (e.g., Strongylura), while a few genera are confined to freshwater rivers and streams, including Belonion, Potamorrhaphis, and Xenentodon. Needlefish closely resemble North American freshwater gars (family Lepisosteidae) in being elongated and having long, narrow jaws filled with sharp teeth, and some species of needlefishes are referred to as gars or garfish despite being only distantly related to the true gars. In fact, the name "garfish" was originally used for the needlefish Belone belone in Europe and only later applied to the North American fishes by European settlers during the 18th century.
Description
Needlefish are slender, ranging from in length. They have a single dorsal fin, placed far back on the body, almost opposite to the anal fin. Their most distinctive feature is their long, narrow beak, which bears multiple sharp teeth. In most species, the upper jaw reaches its full length only in adults, so the juveniles have a half-beak appearance, with an elongated lower jaw, but a much smaller upper one. During this stage of their lifecycle, they eat plankton, switching to fish once the beak fully develops. Needlefish reproduce through mating and laying eggs. The male usually rides the female on the waves as they mate.
Needlefish are most common in the subtropics, but some inhabit temperate waters, as well, particularly during the winter. Belone belone, a common North Atlantic species, often swim in schools alongside tuna. Recently, some small specimens have been seen in the Mediterranean Sea.
Ecology
All needlefish feed primarily on smaller fish. In addition, some species also take krill, swimming crustaceans, small cephalopods and insects.
Needlefish have been documented in taking advantage of Snell's Window when attacking prey; leaping at a shallow angle to ambush schools of small fish. Due to light refraction through water, objects at the edges of the window appear distorted, disrupting the image of the leaping needlefish and allowing it to get within very short distances of its prey.
Danger to humans
Needlefish, like all ray-finned beloniforms, are capable of making short jumps out of the water at up to . Since needlefish swim near the surface, they often leap over the decks of shallow boats rather than going around. This jumping activity is greatly excited by artificial light at night; night fisherman and divers in areas across the Pacific Ocean have been "attacked" by schools of suddenly excited needlefish diving across the water towards the light source at high speed. Their sharp beaks are capable of inflicting deep puncture wounds, often breaking off inside the victim in the process. For many traditional Pacific Islander communities, who primarily fish on reefs from low boats, needlefish represent an even greater risk of injury than sharks.
Occasional deaths and serious injuries have been attributed to needlefish. They include the following documented incidents:
In 1977, a 10-year-old Hawaiian boy, night fishing with his father at Hanamaulu Bay, Kaua'i, was killed when a needlefish jumped from the water and pierced his eye and brain.
In 2007, a 16-year-old Vietnamese boy was stabbed through the heart by the beak of a needlefish while diving for sea cucumbers at night near Halong Bay.
In 2010, a kayaker in Florida was nearly killed when a crocodile needlefish (houndfish, Tylosurus crocodilus) leapt from the water and impaled her in the chest.
In 2012, German kitesurfer Wolfram Reiners was seriously wounded in the foot by a needlefish near the Seychelles.
In May 2013, a kitesurfer in Egypt's Red Sea was speared directly under his knee when a needlefish jumped out of the water.
In October 2013, a Saudi Arabian news website also reported the death of a young Saudi man in Dammam who died of hemorrhaging after being hit by a needlefish on the left side of the neck.
In 2014, a Russian tourist was nearly killed by a needlefish off Nha Trang in Vietnam. The fish bit her neck and left pieces of its teeth inside her spinal cord, paralyzing her.
In early January 2016, a 39-year-old Indonesian woman from Palu, Central Sulawesi, was seriously injured when a half-metre-long needlefish jumped and pierced her just above the right eye. She was swimming in water 80 cm deep in Tanjung Karang, a popular recreational spot in the Donggala Regency, Central Sulawesi. She died a few hours later despite efforts to save her at a local hospital. Shortly after, pictures of her injury spread through instant-messaging applications, while several local news websites also reported the incident, some erroneously attributing the attack to a marlin.
In December 2018, a needlefish struck the neck of a Thai Navy special forces cadet, causing his death.
In April 2024, a 59-year old Filipino fisherman off the coast of Iloilo, Panay Island, Philippines, was pierced in the stomach by a needlefish (locally known as ) jumping out of the water. He was rushed to a nearby hospital but was declared dead on arrival due to hemorrhage.
In October 2024, an Italian surfer died in Indonesia after being impaled in her upper left chest by a needlefish.
In the aquarium
Some species of needlefish inhabit brackish and freshwater environments, and one of the freshwater species, Xenentodon cancila from Southeast Asia, is occasionally kept as an aquarium fish. It is a relatively small species, no more than 40 cm in length when fully grown, but is considered to be a rather delicate fish best suited to advanced aquarists.
Taxonomy
The needlefish family is classified within the order Beloniformes and along with the sauries of the family Scomberesocidae they make up the superfamily Scomberesocoidea and in turn, is one of two superfamilies in the suborder Exocoetoidei which comprises all of the Beloniformes except for the ricefishes of the family Adrianichthydae. Workers have concluded that the genus Belone is the sister group to the sauries and that if this is correct them the Belonidae is only monophyletic if the sauries are included within it.
| Biology and health sciences | Acanthomorpha | Animals |
1226995 | https://en.wikipedia.org/wiki/Granulite | Granulite | Granulites are a class of high-grade metamorphic rocks of the granulite facies that have experienced high-temperature and moderate-pressure metamorphism. They are medium to coarse–grained and mainly composed of feldspars sometimes associated with quartz and anhydrous ferromagnesian minerals, with granoblastic texture and gneissose to massive structure. They are of particular interest to geologists because many granulites represent samples of the deep continental crust. Some granulites experienced decompression from deep in the Earth to shallower crustal levels at high temperature; others cooled while remaining at depth in the Earth.
The minerals present in a granulite will vary depending on the parent rock of the granulite and the temperature and pressure conditions experienced during metamorphism. A common type of granulite found in high-grade metamorphic rocks of the continents contains pyroxene, plagioclase feldspar and accessory garnet, oxides and possibly amphiboles. Both clinopyroxene and orthopyroxene may be present, and in fact, the coexistence of clino- and orthopyroxene in a metabasite (metamorphed basalt) defines the granulite facies.
A granulite may be visually quite distinct with abundant small pink or red pyralspite garnets in a 'granular' holocrystalline matrix. Concentrations of garnets, micas, or amphiboles may form along a linear pattern resembling gneiss or migmatite banding.
Formation
Granulites form at crustal depths, typically during regional metamorphism at high thermal gradients of greater than 30 °C/km. In continental crustal rocks, biotite may break down at high temperatures to form orthopyroxene + potassium feldspar + water, producing a granulite. Other possible minerals formed at dehydration melting conditions include sapphirine, spinel, sillimanite, and osumilite. Some assemblages such as sapphirine + quartz indicate very high temperatures of greater than 900 °C. Some granulites may represent the residues of partial melting at extraction of felsic melts in variable amounts, and in extreme cases represent rocks that all constituent minerals are anhydrous and thus look as if they did not melt at ultrahigh temperature conditions. Therefore, very high temperatures of 900 to 1150 °C are even necessary to produce the granulite-facies mineral assemblages. Such high temperatures at crustal depths only can be delivered by upwelling of the asthenospheric mantle in continental rifting settings, which can cause the regional metamorphism at the high thermal gradients of greater than 30 °C/km.
Granulite facies
The granulite facies is determined by the lower temperature boundary of 700 ± 50 °C and the pressure range of 2–15 kb. The most common mineral assemblage of granulite facies consists of antiperthitic plagioclase, alkali feldspar containing up to 50% albite and Al2O3-rich pyroxenes.
Transition between amphibolite and granulite facies is defined by these reaction isograds:
amphibole → pyroxene + H2O
biotite → K-feldspar + garnet + orthopyroxene + H2O.
Hornblende granulite subfacies is a transitional coexistence region of anhydrous and hydrated ferromagnesian minerals, so the above-mentioned isograds mark the boundary with pyroxene granulite subfacies – facies with completely anhydrous mineral assemblages.
1911 Encyclopædia Britannica definition
Granulite (Latin granulum, "a little grain") is a name used by petrographers to designate two distinct classes of rocks. According to the terminology of the French school it signifies a granite in which both kinds of mica (muscovite and biotite) occur, and corresponds to the German Granit, or to the English muscovite biotite granite. This application has not been accepted generally. [This granitic meaning of granulite is now obsolete.] To the German petrologists granulite means a more or less banded fine-grained metamorphic rock, consisting mainly of quartz and feldspar in very small irregular crystals and usually also containing a fair number of minute, rounded, pale-red garnets. Among English and American geologists the term is generally employed in this sense.
The granulites are very closely allied to the gneisses, as they consist of nearly the same minerals, but they are finer-grained, have usually less perfect foliation, are more frequently garnetiferous, and have some special features of microscopic structure. In the rocks of this group the minerals, as seen in a microscopic slide, occur as small rounded grains forming a closely fitted mosaic. The individual crystals never have perfect form, and indeed traces of it are rare. In some granulites they interlock, with irregular borders; in others they have been drawn out and flattened into tapering lenticles by crushing. In most cases they are somewhat rounded with smaller grains between the larger. This is especially true of the quartz and feldspar which are the predominant minerals; mica always appears as flat scales (irregular or rounded but not hexagonal). Both muscovite and biotite may be present and vary considerably in abundance; very commonly they have their flat sides parallel and give the rock a rudimentary schistosity, and they may be aggregated into bands in which case the granulites are indistinguishable from certain varieties of gneiss. The garnets are very generally larger than the above-mentioned ingredients, and easily visible with the eye as pink spots on the broken surfaces of the rock. They usually are filled with enclosed grains of the other minerals.
The feldspar of the granulites is mostly orthoclase or cryptoperthite; microcline, oligoclase and albite are also common. Basic feldspars occur only rarely. Among accessory minerals, in addition to apatite, zircon, and iron oxides, the following may be mentioned: hornblende (not common), riebeckite (rare), epidote and zoisite, calcite, sphene, andalusite, sillimanite, kyanite, hercynite (a green spinel), rutile, orthite and tourmaline. Though occasionally we may find larger grains of feldspar, quartz or epidote, it is more characteristic of these rocks that all the minerals are in small, nearly uniform, imperfectly shaped individuals.
On account of the minuteness with which it has been described and the important controversies on points of theoretical geology which have arisen regarding it, the granulite district of Saxony (in the area of Rosswein and Penig) in Germany may be considered the typical region for rocks of this group. It should be remembered that though granulites are probably the commonest rocks of this country, they are mingled with granites, gneisses, gabbros, amphibolites, mica schists and many other petrographical types. All of these rocks show more-or-less metamorphism either of a thermal character or due to pressure and crushing. The granites pass into gneiss and granulite; the gabbros into flaser gabbro and amphibolite; the slates often contain andalusite or chiastolite, and show transitions to mica schists. At one time these rocks were regarded as Archean gneisses of a special type. Johannes Georg Lehmann propounded the hypothesis that their present state was due principally to crushing acting on them in a solid condition, grinding them down and breaking up their minerals, while the pressure to which they were subjected welded them together into coherent rock. It is now believed, however, that they are comparatively recent and include sedimentary rocks, partly of Palaeozoic age, and intrusive masses which may be nearly massive or may have gneissose, flaser or granulitic structures. These have been developed largely by the injection of semi-consolidated highly viscous intrusions, and the varieties of texture are original or were produced very shortly after the crystallization of the rocks. Meanwhile, however, Lehmanns advocacy of post-consolidation crushing as a factor in the development of granulites has been so successful that the terms granulitization and granulitic structures are widely employed to indicate the results of dynamometamorphism acting on rocks at a period long after their solidification.
The Saxon granulites are apparently for the most part igneous and correspond in composition to granites and porphyries. There are, however, many granulites which undoubtedly were originally sediments (arkoses, grits and sandstones). A large part of the highlands of Scotland consists of paragranulites of this kind, which have received the group name of Moine gneisses.
Along with the typical acid granulites above described, in Saxony, India, Scotland and other countries there occur dark-colored basic granulites (trap granulites). These are fine-grained rocks, not usually banded, nearly black in color with small red spots of garnet. Their essential minerals are pyroxene, plagioclase and garnet: chemically they resemble the gabbros. Green augite and hypersthene form a considerable part of these rocks, they may contain also biotite, hornblende and quartz. Around the garnets there is often a radial grouping of small grains of pyroxene and hornblende in a clear matrix of feldspar: these centric structures are frequent in granulites. The rocks of this group accompany gabbro and serpentine, but the exact conditions under which they are formed and the significance of their structures is not very clearly understood.
| Physical sciences | Metamorphic rocks | Earth science |
1227081 | https://en.wikipedia.org/wiki/Phosphorite | Phosphorite | Phosphorite, phosphate rock or rock phosphate is a non-detrital sedimentary rock that contains high amounts of phosphate minerals. The phosphate content of phosphorite (or grade of phosphate rock) varies greatly, from 4% to 20% phosphorus pentoxide (P2O5). Marketed phosphate rock is enriched ("beneficiated") to at least 28%, often more than 30% P2O5. This occurs through washing, screening, de-liming, magnetic separation or flotation. By comparison, the average phosphorus content of sedimentary rocks is less than 0.2%.
The phosphate is present as fluorapatite Ca5(PO4)3F typically in cryptocrystalline masses (grain sizes < 1 μm) referred to as collophane-sedimentary apatite deposits of uncertain origin. It is also present as hydroxyapatite Ca5(PO4)3OH or Ca10(PO4)6(OH)2, which is often dissolved from vertebrate bones and teeth, whereas fluorapatite can originate from hydrothermal veins. Other sources also include chemically dissolved phosphate minerals from igneous and metamorphic rocks. Phosphorite deposits often occur in extensive layers, which cumulatively cover tens of thousands of square kilometres of the Earth's crust.
Limestones and mudstones are common phosphate-bearing rocks. Phosphate-rich sedimentary rocks can occur in dark brown to black beds, ranging from centimeter-sized laminae to beds that are several meters thick. Although these thick beds can exist, they are rarely only composed of phosphatic sedimentary rocks. Phosphatic sedimentary rocks are commonly accompanied by or interbedded with shales, cherts, limestone, dolomites and sometimes sandstone. These layers contain the same textures and structures as fine-grained limestones and may represent diagenetic replacements of carbonate minerals by phosphates. They also can be composed of peloids, ooids, fossils, and clasts that are made up of apatite. There are some phosphorites that are very small and have no distinctive granular textures. This means that their textures are similar to that of collophane, or fine micrite-like texture. Phosphatic grains may be accompanied by organic matter, clay minerals, silt-sized detrital grains, and pyrite. Peloidal or pelletal phosphorites occur normally; whereas oolitic phosphorites are not common.
Phosphorites are known from Proterozoic banded iron formations in Australia, but are more common from Paleozoic and Cenozoic sediments. The Permian Phosphoria Formation of the western United States represents some 15 million years of sedimentation. It reaches a thickness of 420 metres and covers an area of 350,000 km2. Commercially mined phosphorites occur in France, Belgium, Spain, Morocco, Tunisia, Saudi Arabia and Algeria. In the United States phosphorites have been mined in Florida, Tennessee, Wyoming, Utah, Idaho and Kansas.
Classification of phosphatic sedimentary rocks
(1) Pristine: Phosphates that are in pristine conditions have not undergone bioturbation. In other words, the word pristine is used when phosphatic sediment, phosphatized stromatolites and phosphate hardgrounds have not been disturbed.
(2) Condensed: Phosphatic particles, laminae and beds are considered condensed when they have been concentrated. This is helped by the extracting and reworking processes of phosphatic particles or bioturbation.
(3) Allochthonous: Phosphatic particles that were moved by turbulent or gravity-driven flows and deposited by these flows.
Phosphorus cycle, formation and accumulation
The heaviest accumulation of phosphorus is mainly on the ocean floor. Phosphorus accumulation occurs from atmospheric precipitation, dust, glacial runoff, cosmic activity, underground hydrothermal volcanic activity, and deposition of organic material. The primary inflow of dissolved phosphorus is from continental weathering, brought out by rivers to the ocean. It is then processed by both micro- and macro-organisms. Diatomaceous plankton, phytoplankton, and zooplankton process and dissolve phosphorus in the water. The bones and teeth of certain fish (e.g. anchovies) absorb phosphorus and are later deposited and buried in the marine sediment.
Depending on the pH and salinity levels of the ocean water, organic matter will decay, releasing phosphorus from sediment in shallow basins. Bacteria and enzymes dissolve organic matter on the water–bottom interface, thus returning phosphorus to the beginning of its biogenic cycle. Mineralization of organic matter can also cause the release of phosphorus back into the ocean water.
Depositional environments
Phosphates are known to be deposited in a wide range of depositional environments. Normally phosphates are deposited in very shallow, near-shore marine or low energy environments. This includes environments such as supratidal zones, littoral or intertidal zones, and most importantly estuarine. Currently, areas of oceanic upwelling cause the formation of phosphates. This is because of the constant stream of phosphate brought from the large, deep ocean reservoir (see below). This cycle allows continuous growth of organisms.
Supratidal zones: Supratidal environments are part of the tidal flat system where the presence of strong wave activity is non-existent. Tidal flat systems are created along open coasts and relatively low wave energy environments. They can also develop on high energy coasts behind barrier islands where they are sheltered from the high energy wave action. Within the tidal flat system, the supratidal zone lies in a very high tide level. However, it can be flooded by extreme tides and cut across by tidal channels. This is also subaerially exposed, but is flooded twice a month by spring tides.
Littoral environments/intertidal zones: Intertidal zones are also part of the tidal flat system. The intertidal zone is located within the mean high and low tide levels. It is subject to tidal shifts, which means that it is subaerially exposed once or twice a day. It is not exposed long enough to support the growth of most vegetation. The zone contains both suspension sedimentation and bed load.
Estuarine environments: Estuarine environments, or estuaries, are located at the lower parts of rivers that flow into the open sea. Since they are in the seaward section of the drowned valley system they receive sediment from both marine and fluvial sources. These contain facies that are affected by tide and wave fluvial processes. An estuary is considered to stretch from the landward limit of tidal facies to the seaward limit of coastal facies. Phosphorites are often deposited in fjords within estuarine environments. These are estuaries with shallow sill constrictions. During Holocene sea-level rise, fjord estuaries formed by drowning of glacially-eroded U-shaped valleys.
The most common occurrence of phosphorites is related to strong marine upwelling of sediments. Upwelling is caused by deep water currents that are brought up to coastal surfaces where a large deposition of phosphorites may occur. This type of environment is the main reason why phosphorites are commonly associated with silica and chert. Estuaries are also known as a phosphorus “trap”. This is because coastal estuaries contain a high productivity of phosphorus from marsh grass and benthic algae which allow an equilibrium exchange between living and dead organisms.
Types of phosphorite deposition
Phosphate nodules: These are spherical concentrations that are randomly distributed along the floor of continental shelves. Most phosphorite grains are sand size although particles greater than 2 mm may be present. These larger grains, referred to as nodules, can range up to several tens of centimeters in size. Phosphate nodules are known to occur in significant quantities offshore northern Chile.
Bioclastic phosphates or bone beds: Bone beds are bedded phosphate deposits that contain concentrations of small skeletal particles and coprolites. Some also contain invertebrate fossils like brachiopods and become more enriched in P2O5 after diagenetic processes have occurred. Bioclastic phosphates can also be cemented by phosphate minerals.
Phosphatization: Phosphatization is a type of rare diagenetic processes. It occurs when fluids that are rich in phosphate are leached from guano. These are then concentrated and reprecipitated in limestone. Phosphatized fossils or fragments of original phosphatic shells are important components within some these deposits.
Tectonic and oceanographic settings of marine phosphorites
Epeiric sea phosphorites: Epeiric sea phosphorites are within marine shelf environments. These are in a broad and shallow cratonic setting. This is where granular phosphorites, phosphorite hardgrounds, and nodules occur.
Continental margin phosphorites: Convergent, passive, upwelling, non-upwelling. This environment accumulates phosphorites in the form of hardgrounds, nodules and granular beds. These accumulate by carbonate fluorapatite precipitation during early diagenesis in the upper few tens of centimeters of sediment. There are two different environmental conditions in which phosphorites are produced within continental margins. Continental margins can consist of organic rich sedimentation, strong coastal upwelling, and pronounced low oxygen zones. They can also form in conditions such as oxygen rich bottom waters and organic poor sediments.
Seamount phosphorites: These are phosphorites that occur in seamounts, guyots, or flat topped seamounts, seamount ridges. These phosphorites are produced in association with iron and magnesium bearing crusts. In this setting the productivity of phosphorus is recycled within an iron oxidation reduction phosphorus cycle. This cycle can also form glauconite which is normally associated with modern and ancient phosphorites.
Insular phosphorites: Insular phosphorites are located in carbonate islands, plateaus, coral island consisting of a reef surrounding a lagoon or, atoll lagoon, marine lakes. The phosphorite here originates from guano. Replacement of deep sea sediments precipitates, that has been formed in place on the ocean floor.
Production and use
Production
Deposits which contain phosphate in quantity and concentration which are economic to mine as ore for their phosphate content are not particularly common. The two main sources for phosphate are guano, formed from bird or bat droppings, and rocks containing concentrations of the calcium phosphate mineral, apatite.
, the US is the world's leading producer and exporter of phosphate fertilizers, accounting for about 37% of world P2O5 exports. , the world's total economic demonstrated resource of rock phosphate is 70 gigatonnes, which occurs principally as sedimentary marine phosphorites.
, China, the United States and Morocco are the world's largest miners of phosphate rock, with a production of 77 megatonnes, 29.4 Mt and 26.8 Mt (including 2.5 Mt in the Sahara of Morocco) respectively in 2012 while global production reached 195 Mt. It is thought that in India there are almost 260 million tons of rock phosphate. Other countries with significant production include Brazil, Russia, Jordan and Tunisia. Historically, large amounts of phosphates were obtained from deposits on small islands such as Christmas Island and Nauru, but these sources are now largely depleted.
Phosphate ore is mined and beneficiated into rock phosphate. Beneficiation of phosphate ore is a process which includes washing, flotation and calcining. Froth flotation is used to concentrate the mined ore to rock phosphate. The mined ore is crushed and washed, creating a slurry, this ore slurry is then treated with fatty acids to cause calcium phosphate to become hydrophobic.
This rock phosphate is then either solubilized to produce wet-process phosphoric acid, or smelted to produce elemental phosphorus. Phosphoric acid is reacted with phosphate rock to produce the fertilizer triple superphosphate or with anhydrous ammonia to produce the ammonium phosphate fertilizers. Elemental phosphorus is the base for furnace-grade phosphoric acid, phosphorus pentasulfide, phosphorus pentoxide, and phosphorus trichloride.
Uses
Approximately 90% of rock phosphate production is used for fertilizer and animal feed supplements and the balance for industrial chemicals. In addition, phosphorus from rock phosphate is also used in food preservatives, baking flour, pharmaceuticals, anticorrosion agents, cosmetics, fungicides, insecticides, detergents, ceramics, water treatment and metallurgy.
For use in the chemical fertilizer industry, beneficiated rock phosphate must be concentrated to levels of at least 28% phosphorus pentoxide (P2O5), although most marketed grades of phosphate rock are 30% or more.
It must also have reasonable amounts of calcium carbonate (5%), and <4% combined iron and aluminium oxides. Worldwide, the resources of high-grade ore are declining, and use of lower grade ore may become more attractive.
Beneficiated rock phosphate is also marketed and accepted as an "organic" alternative to "chemical" phosphate fertilizer which has been further concentrated from it, because it is perceived as being more "natural". According to a report for the FAO, it can be more sustainable to apply rock phosphate as a fertilizer in certain soil types and countries, although it has many drawbacks. According to the report it may have higher sustainability compared to more concentrated fertilizers because of reduced manufacturing costs and the possibility of local procurement of the refined ore.
Rare earth elements are being found within phosphorites. With increasing demand from modern technology a different method of finding rare earth elements, independent of China, is becoming increasingly important. With yields greater than those from deposits in China, phosphorites offer a new resource located within the U.S. that would likely lead to independence from influence of countries outside of the U.S.
| Physical sciences | Sedimentary rocks | Earth science |
1227284 | https://en.wikipedia.org/wiki/Pulsar | Pulsar | A pulsar (from puls(ating st)ar, on the model of quasar) is a highly magnetized rotating neutron star that emits beams of electromagnetic radiation out of its magnetic poles. This radiation can be observed only when a beam of emission is pointing toward Earth (similar to the way a lighthouse can be seen only when the light is pointed in the direction of an observer), and is responsible for the pulsed appearance of emission. Neutron stars are very dense and have short, regular rotational periods. This produces a very precise interval between pulses that ranges from milliseconds to seconds for an individual pulsar. Pulsars are one of the candidates for the source of ultra-high-energy cosmic rays (see also centrifugal mechanism of acceleration.)
Pulsars’ highly regular pulses make them very useful tools for astronomers. For example, observations of a pulsar in a binary neutron star system were used to indirectly confirm the existence of gravitational radiation. The first extrasolar planets were discovered in 1992 around a pulsar, specifically PSR B1257+12. In 1983, certain types of pulsars were detected that, at that time, exceeded the accuracy of atomic clocks in keeping time.
History of observation
Discovery
Signals from the first discovered pulsar were initially observed by Jocelyn Bell while analyzing data recorded on August 6, 1967, from a newly commissioned radio telescope that she helped build. Initially dismissed as radio interference by her supervisor and developer of the telescope, Antony Hewish, the fact that the signals always appeared at the same declination and right ascension soon ruled out a terrestrial source. On November 28, 1967, Bell and Hewish using a fast strip chart recorder resolved the signals as a series of pulses, evenly spaced every 1.337 seconds. No astronomical object of this nature had ever been observed before. On December 21, Bell discovered a second pulsar, quashing speculation that these might be signals beamed at earth from an extraterrestrial intelligence.
When observations with another telescope confirmed the emission, it eliminated any sort of instrumental effects. At this point, Bell said of herself and Hewish that "we did not really believe that we had picked up signals from another civilization, but obviously the idea had crossed our minds and we had no proof that it was an entirely natural radio emission. It is an interesting problem—if one thinks one may have detected life elsewhere in the universe, how does one announce the results responsibly?" Even so, they nicknamed the signal LGM-1, for "little green men" (a playful name for intelligent beings of extraterrestrial origin).
It was not until a second pulsating source was discovered in a different part of the sky that the "LGM hypothesis" was entirely abandoned. Their pulsar was later dubbed CP 1919, and is now known by a number of designators including PSR B1919+21 and PSR J1921+2153. Although CP 1919 emits in radio wavelengths, pulsars have subsequently been found to emit in visible light, X-ray, and gamma ray wavelengths.
The word "pulsar" first appeared in print in 1968:
The existence of neutron stars was first proposed by Walter Baade and Fritz Zwicky in 1934, when they argued that a small, dense star consisting primarily of neutrons would result from a supernova. Based on the idea of magnetic flux conservation from magnetic main sequence stars, Lodewijk Woltjer proposed in 1964 that such neutron stars might contain magnetic fields as large as 1014 to 1016 gauss (=1010 to 1012 tesla). In 1967, shortly before the discovery of pulsars, Franco Pacini suggested that a rotating neutron star with a magnetic field would emit radiation, and even noted that such energy could be pumped into a supernova remnant around a neutron star, such as the Crab Nebula. After the discovery of the first pulsar, Thomas Gold independently suggested a rotating neutron star model similar to that of Pacini, and explicitly argued that this model could explain the pulsed radiation observed by Bell Burnell and Hewish.
In 1968, Richard V. E. Lovelace with collaborators discovered period ms of the Crab Nebula pulsar using Arecibo Observatory.
The discovery of the Crab pulsar provided confirmation of the rotating neutron star model of pulsars. The Crab pulsar 33-millisecond pulse period was too short to be consistent with other proposed models for pulsar emission. Moreover, the Crab pulsar is so named because it is located at the center of the Crab Nebula, consistent with the 1933 prediction of Baade and Zwicky.
In 1974, Antony Hewish and Martin Ryle, who had developed revolutionary radio telescopes, became the first astronomers to be awarded the Nobel Prize in Physics, with the Royal Swedish Academy of Sciences noting that Hewish played a "decisive role in the discovery of pulsars". Considerable controversy is associated with the fact that Hewish was awarded the prize while Bell, who made the initial discovery while she was his PhD student, was not. Bell claims no bitterness upon this point, supporting the decision of the Nobel prize committee.
Milestones
In 1974, Joseph Hooton Taylor, Jr. and Russell Hulse discovered for the first time a pulsar in a binary system, PSR B1913+16. This pulsar orbits another neutron star with an orbital period of just eight hours. Einstein's theory of general relativity predicts that this system should emit strong gravitational radiation, causing the orbit to continually contract as it loses orbital energy. Observations of the pulsar soon confirmed this prediction, providing the first ever evidence of the existence of gravitational waves. As of 2010, observations of this pulsar continue to agree with general relativity. In 1993, the Nobel Prize in Physics was awarded to Taylor and Hulse for the discovery of this pulsar.
In 1982, Don Backer led a group that discovered PSR B1937+21, a pulsar with a rotation period of just 1.6 milliseconds (38,500 rpm). Observations soon revealed that its magnetic field was much weaker than ordinary pulsars, while further discoveries cemented the idea that a new class of object, the "millisecond pulsars" (MSPs) had been found. MSPs are believed to be the end product of X-ray binaries. Owing to their extraordinarily rapid and stable rotation, MSPs can be used by astronomers as clocks rivaling the stability of the best atomic clocks on Earth. Factors affecting the arrival time of pulses at Earth by more than a few hundred nanoseconds can be easily detected and used to make precise measurements. Physical parameters accessible through pulsar timing include the 3D position of the pulsar, its proper motion, the electron content of the interstellar medium along the propagation path, the orbital parameters of any binary companion, the pulsar rotation period and its evolution with time. (These are computed from the raw timing data by Tempo, a computer program specialized for this task.) After these factors have been taken into account, deviations between the observed arrival times and predictions made using these parameters can be found and attributed to one of three possibilities: intrinsic variations in the spin period of the pulsar, errors in the realization of Terrestrial Time against which arrival times were measured, or the presence of background gravitational waves. Scientists are currently attempting to resolve these possibilities by comparing the deviations seen between several different pulsars, forming what is known as a pulsar timing array. The goal of these efforts is to develop a pulsar-based time standard precise enough to make the first ever direct detection of gravitational waves. In 2006, a team of astronomers at LANL proposed a model to predict the likely date of pulsar glitches with observational data from the Rossi X-ray Timing Explorer. They used observations of the pulsar PSR J0537−6910, that is known to be a quasi-periodic glitching pulsar. However, no general scheme for glitch forecast is known to date.
In 1992, Aleksander Wolszczan discovered the first extrasolar planets around PSR B1257+12. This discovery presented important evidence concerning the widespread existence of planets outside the Solar System, although it is very unlikely that any life form could survive in the environment of intense radiation near a pulsar.
Pulsar-like white dwarfs
White dwarfs can also act as pulsars. Because the moment of inertia of a white dwarf is much higher than that of a neutron star, the white-dwarf pulsars rotate once every several minutes, far slower than neutron-star pulsars.
By 2024, three pulsar-like white dwarfs have been identified.
In 1998, Nazar Ikhsanov showed that a white dwarf in the binary system AE Aquarii acts like a radio pulsar. The confirmation of the pulsar-like properties of the white dwarf in AE Aquarii was provided in 2008 by a discovery of X-ray pulsations, which showed that this white dwarf acts not only as a radio pulsar, but also as an X-ray pulsar.
In 2016, a white dwarf in the binary system AR Scorpii was identified as a pulsar (it is often mistakenly called the first discovered pulsar-like white dwarf). The system displays strong pulsations from ultraviolet to radio wavelengths, powered by the spin-down of the strongly magnetized white dwarf.
In 2023, it was suggested that the white dwarf eRASSU J191213.9−441044 acts as a pulsar both in radio and X-rays.
There is an alternative tentative explanation of the pulsar-like properties of these white dwarfs. In 2019, the properties of pulsars have been explained using a numerical magnetohydrodynamic model explaining was developed at Cornell University. According to this model, AE Aqr is an intermediate polar-type star, where the magnetic field is relatively weak and an accretion disc may form around the white dwarf. The star is in the propeller regime, and many of its observational properties are determined by the disc-magnetosphere interaction. A similar model for eRASSU J191213.9−441044 is supported by the results of its observations at ultraviolet wave lengths, which showed that its magnetic field strength does not exceed 50 MG.
Nomenclature
Initially pulsars were named with letters of the discovering observatory followed by their right ascension (e.g. CP 1919). As more pulsars were discovered, the letter code became unwieldy, and so the convention then arose of using the letters PSR (Pulsating Source of Radio) followed by the pulsar's right ascension and degrees of declination (e.g. PSR 0531+21) and sometimes declination to a tenth of a degree (e.g. PSR 1913+16.7). Pulsars appearing very close together sometimes have letters appended (e.g. PSR 0021−72C and PSR 0021−72D).
The modern convention prefixes the older numbers with a B (e.g. PSR B1919+21), with the B meaning the coordinates are for the 1950.0 epoch. All new pulsars have a J indicating 2000.0 coordinates and also have declination including minutes (e.g. PSR J1921+2153). Pulsars that were discovered before 1993 tend to retain their B names rather than use their J names (e.g. PSR J1921+2153 is more commonly known as PSR B1919+21). Recently discovered pulsars only have a J name (e.g. PSR J0437−4715). All pulsars have a J name that provides more precise coordinates of its location in the sky.
Formation, mechanism, turn off
The events leading to the formation of a pulsar begin when the core of a massive star is compressed during a supernova, which collapses into a neutron star. The neutron star retains most of its angular momentum, and since it has only a tiny fraction of its progenitor's radius, it is formed with very high rotation speed. A beam of radiation is emitted along the magnetic axis of the pulsar, which spins along with the rotation of the neutron star. The magnetic axis of the pulsar determines the direction of the electromagnetic beam, with the magnetic axis not necessarily being the same as its rotational axis. This misalignment causes the beam to be seen once for every rotation of the neutron star, which leads to the "pulsed" nature of its appearance.
In rotation-powered pulsars, the beam is the result of the rotational energy of the neutron star, which generates an electrical field and very strong magnetic field, resulting in the acceleration of protons and electrons on the star surface and the creation of an electromagnetic beam emanating from the poles of the magnetic field. Observations by NICER of PSR J0030+0451 indicate that both beams originate from hotspots located on the south pole and that there may be more than two such hotspots on that star. This rotation slows down over time as electromagnetic power is emitted. When a pulsar's spin period slows down sufficiently, the radio pulsar mechanism is believed to turn off (the so-called "death line"). This turn-off seems to take place after about 10–100 million years, which means of all the neutron stars born in the 13.6-billion-year age of the universe, around 99% no longer pulsate.
Though the general picture of pulsars as rapidly rotating neutron stars is widely accepted, Werner Becker of the Max Planck Institute for Extraterrestrial Physics said in 2006, "The theory of how pulsars emit their radiation is still in its infancy, even after nearly forty years of work."
Categories
Three distinct classes of pulsars are currently known to astronomers, according to the source of the power of the electromagnetic radiation:
rotation-powered pulsars, where the loss of rotational energy of the star provides the power,
accretion-powered pulsars (accounting for most but not all X-ray pulsars), where the gravitational potential energy of accreted matter is the power source (producing X-rays that are observable from the Earth),
magnetars, where the decay of an extremely strong magnetic field provides the electromagnetic power.
Although all three classes of objects are neutron stars, their observable behavior and the underlying physics are quite different. There are, however, some connections. For example, X-ray pulsars are probably old rotationally-powered pulsars that have already lost most of their energy, and have only become visible again after their binary companions had expanded and begun transferring matter on to the neutron star.
The process of accretion can, in turn, transfer enough angular momentum to the neutron star to "recycle" it as a rotation-powered millisecond pulsar. As this matter lands on the neutron star, it is thought to "bury" the magnetic field of the neutron star (although the details are unclear), leaving millisecond pulsars with magnetic fields 1000–10,000 times weaker than average pulsars. This low magnetic field is less effective at slowing the pulsar's rotation, so millisecond pulsars live for billions of years, making them the oldest known pulsars. Millisecond pulsars are seen in globular clusters, which stopped forming neutron stars billions of years ago.
Of interest to the study of the state of the matter in a neutron star are the glitches observed in the rotation velocity of the neutron star. This velocity decreases slowly but steadily, except for an occasional sudden variation known as "glitch". One model put forward to explain these glitches is that they are the result of "starquakes" that adjust the crust of the neutron star. Models where the glitch is due to a decoupling of the possibly superconducting interior of the star have also been advanced. In both cases, the star's moment of inertia changes, but its angular momentum does not, resulting in a change in rotation rate.
Disrupted recycled pulsar
When two massive stars are born close together from the same cloud of gas, they can form a binary system and orbit each other from birth. If those two stars are at least a few times as massive as the Sun, their lives will both end in supernova explosions. The more massive star explodes first, leaving behind a neutron star. If the explosion does not kick the second star away, the binary system survives. The neutron star can now be visible as a radio pulsar, and it slowly loses energy and spins down. Later, the second star can swell up, allowing the neutron star to suck up its matter. The matter falling onto the neutron star spins it up and reduces its magnetic field.
This is called "recycling" because it returns the neutron star to a quickly-spinning state. Finally, the second star also explodes in a supernova, producing another neutron star. If this second explosion also fails to disrupt the binary, a double neutron star (neutron star binary) is formed. Otherwise, the spun-up neutron star is left with no companion and becomes a "disrupted recycled pulsar", spinning between a few and 50 times per second.
Applications
The discovery of pulsars allowed astronomers to study an object never observed before, the neutron star. This kind of object is the only place where the behavior of matter at nuclear density can be observed (though not directly). Also, millisecond pulsars have allowed a test of general relativity in conditions of an intense gravitational field.
Maps
Pulsar maps have been included on the two Pioneer plaques as well as the Voyager Golden Record. They show the position of the Sun, relative to 14 pulsars, which are identified by the unique timing of their electromagnetic pulses, so that Earth's position both in space and time can be calculated by potential extraterrestrial intelligence. Because pulsars are emitting very regular pulses of radio waves, its radio transmissions do not require daily corrections. Moreover, pulsar positioning could create a spacecraft navigation system independently, or be used in conjunction with satellite navigation.
Pulsar navigation
X-ray pulsar-based navigation and timing (XNAV) or simply pulsar navigation is a navigation technique whereby the periodic X-ray signals emitted from pulsars are used to determine the location of a vehicle, such as a spacecraft in deep space. A vehicle using XNAV would compare received X-ray signals with a database of known pulsar frequencies and locations. Similar to GPS, this comparison would allow the vehicle to calculate its position accurately (±5 km). The advantage of using X-ray signals over radio waves is that X-ray telescopes can be made smaller and lighter. Experimental demonstrations have been reported in 2018.
Precise clocks
Generally, the regularity of pulsar emission does not rival the stability of atomic clocks. They can still be used as external reference. For example, J0437−4715 has a period of s with an error of .
This stability allows millisecond pulsars to be used in establishing ephemeris time
or in building pulsar clocks.
Timing noise is the name for rotational irregularities observed in all pulsars. This timing noise is observable as random wandering in the pulse frequency or phase. It is unknown whether timing noise is related to pulsar glitches. According to a study published in 2023, the timing noise observed in pulsars is believed to be caused by background gravitational waves. Alternatively, it may be caused by stochastic fluctuations in both the internal (related to the presence of superfluids or turbulence) and external (due to magnetospheric activity) torques in a pulsar.
Probes of the interstellar medium
The radiation from pulsars passes through the interstellar medium (ISM) before reaching Earth. Free electrons in the warm (8000 K), ionized component of the ISM and H II regions affect the radiation in two primary ways. The resulting changes to the pulsar's radiation provide an important probe of the ISM itself.
Because of the dispersive nature of the interstellar plasma, lower-frequency radio waves travel through the medium slower than higher-frequency radio waves. The resulting delay in the arrival of pulses at a range of frequencies is directly measurable as the dispersion measure of the pulsar. The dispersion measure is the total column density of free electrons between the observer and the pulsar:
where is the distance from the pulsar to the observer, and is the electron density of the ISM. The dispersion measure is used to construct models of the free electron distribution in the Milky Way.
Additionally, density inhomogeneities in the ISM cause scattering of the radio waves from the pulsar. The resulting scintillation of the radio waves—the same effect as the twinkling of a star in visible light due to density variations in the Earth's atmosphere—can be used to reconstruct information about the small scale variations in the ISM. Due to the high velocity (up to several hundred km/s) of many pulsars, a single pulsar scans the ISM rapidly, which results in changing scintillation patterns over timescales of a few minutes. The exact cause of these density inhomogeneities remains an open question, with possible explanations ranging from turbulence to current sheets.
Probes of space-time
Pulsars orbiting within the curved space-time around Sgr A*, the supermassive black hole at the center of the Milky Way, could serve as probes of gravity in the strong-field regime. Arrival times of the pulses would be affected by special- and general-relativistic Doppler shifts and by the complicated paths that the radio waves would travel through the strongly curved space-time around the black hole. In order for the effects of general relativity to be measurable with current instruments, pulsars with orbital periods less than about 10 years would need to be discovered; such pulsars would orbit at distances inside 0.01 pc from Sgr A*. Searches are currently underway; at present, five pulsars are known to lie within 100 pc from Sgr A*.
Gravitational wave detectors
There are four consortia around the world which use pulsars to search for gravitational waves: the European Pulsar Timing Array (EPTA) in Europe, the Parkes Pulsar Timing Array (PPTA) in Australia, the North American Nanohertz Observatory for Gravitational Waves (NANOGrav) in Canada and the US, and the Indian Pulsar Timing Array (InPTA) in India. Together, the consortia form the International Pulsar Timing Array (IPTA). The pulses from Millisecond Pulsars (MSPs) are used as a system of galactic clocks. Disturbances in the clocks will be measurable at Earth. A disturbance from a passing gravitational wave will have a particular signature across the ensemble of pulsars, and will be thus detected.
Significant pulsars
The pulsars listed here were either the first discovered of its type, or represent an extreme of some type among the known pulsar population, such as having the shortest measured period.
The first radio pulsar "CP 1919" (now known as PSR B1919+21), with a pulse period of 1.337 seconds and a pulse width of 0.04-second, was discovered in 1967.
The first binary pulsar, PSR 1913+16, whose orbit is decaying due to the emission of gravitational radiation at the exact rate predicted by general relativity.
The brightest radio pulsar, the Vela Pulsar.
The first millisecond pulsar, PSR B1937+21
The brightest millisecond pulsar, PSR J0437−4715
The first X-ray pulsar, Cen X-3
The first accreting millisecond X-ray pulsar, SAX J1808.4−3658
The first pulsar with planets, PSR B1257+12
The first pulsar observed to have been affected by asteroids: PSR J0738−4042
The first double pulsar binary system, PSR J0737−3039
The shortest period pulsar, PSR J1748−2446ad, with a period of ~0.0014 seconds or ~1.4 milliseconds (716 times a second).
The longest period neutron star pulsar, PSR J0901-4046, with a period of 75.9 seconds.
The longest period pulsar, at 118.2 seconds, as well as one of the only known two white dwarf pulsars, AR Scorpii.
The first white dwarf pulsar AE Aquarii.
The pulsar with the most stable period, PSR J0437−4715
The first millisecond pulsar with 2 stellar mass companions, PSR J0337+1715
PSR J1841−0500, stopped pulsing for 580 days. One of only two pulsars known to have stopped pulsing for more than a few minutes.
PSR B1931+24, has a cycle. It pulses for about a week and stops pulsing for about a month. One of only two pulsars known to have stopped pulsing for more than a few minutes.
Swift J0243.6+6124 most magnetic pulsar with .
PSR J0952-0607 heaviest pulsar with .
PSR J1903+0327, a ~2.15 ms pulsar discovered to be in a highly eccentric binary star system with a Sun-like star.
PSR J2007+2722, a 40.8-hertz 'recycled' isolated pulsar was the first pulsar found by volunteers on data taken in February 2007 and analyzed by distributed computing project Einstein@Home.
PSR J1311–3430, the first millisecond pulsar discovered via gamma-ray pulsations and part of a binary system with the shortest orbital period.
Pulsars in The Quran
In Surah At-Tariq (The Nightcomer), the Quran mentions a celestial body described as "the star of piercing brightness" (Quran 86:3). This description has been interpreted by some to refer to pulsars, highly magnetized, rotating neutron stars that emit beams of electromagnetic radiation. The rhythmic pulses of these stars are often likened to a "knocking" sound, aligning with the term "Tariq," which can mean "knocker" or "visitor" in Arabic. This interpretation suggests that the Quran referenced pulsars over 1,400 years ago, long before their discovery by modern science.
However, this interpretation is debated among scholars. Some argue that the term "Tariq" refers to a bright star visible at night, such as Venus, and does not specifically denote pulsars. They emphasize that the Quran's description aligns with observable celestial phenomena known at the time.
Regardless of the debate, the mention of a "piercing star" in the Quran highlights the text's poetic and metaphorical language, which has inspired various interpretations over the centuries.
Gallery
| Physical sciences | Stellar astronomy | null |
1227385 | https://en.wikipedia.org/wiki/Buddhist%20calendar | Buddhist calendar | The Buddhist calendar is a set of lunisolar calendars primarily used in Tibet, Cambodia, Laos, Myanmar, Bangladesh, India, Sri Lanka, Thailand and Vietnam as well as in Malaysia and Singapore and by Chinese populations for religious or official occasions. While the calendars share a common lineage, they also have minor but important variations such as intercalation schedules, month names and numbering, use of cycles, etc. In Thailand, the name Buddhist Era is a year numbering system shared by the traditional Thai lunar calendar and by the Thai solar calendar.
The Southeast Asian lunisolar calendars are largely based on an older version of the Hindu calendar, which uses the sidereal year as the solar year. One major difference is that the Southeast Asian systems, unlike their Indian cousins, do not use apparent reckoning to stay in sync with the sidereal year. Instead, they employ their versions of the Metonic cycle. However, since the Metonic cycle is not very accurate for sidereal years, the Southeast Asian calendar is slowly drifting out of sync with the sidereal, approximately one day every 100 years. Yet no coordinated structural reforms of the lunisolar calendar have been undertaken.
Today, the traditional Buddhist lunisolar calendar is used mainly for Theravada Buddhist festivals. The Thai Buddhist Era, a renumbered Gregorian calendar, is the official calendar in Thailand.
Structure
The calculation methodology of the current versions of Southeast Asian Buddhist calendars is largely based on that of the Burmese calendar, which was in use in various Southeast Asian kingdoms down to the 19th century under the names of Chula Sakarat and Jolak Sakaraj. The Burmese calendar in turn was based on the "original" Surya Siddhanta system of ancient India (believed to be Ardharatrika school). One key difference with Indian systems is that the Burmese system has followed a variation of the Metonic cycle. It is unclear from where, when or how the Metonic system was introduced; hypotheses range from China to Europe. The Burmese system, and indeed the Southeast Asian systems, thus use a "strange" combination of sidereal years from the Indian calendar in combination with the Metonic cycle better for tropical years.
Epochal date
In all Theravada traditions,the calendar's epochal year 0 date was the day in which the Buddha attained parinibbāna. However, not all traditions agree on when it actually took place. In Burmese Buddhist tradition, it was 13 May 554 BCE (Tuesday, Full moon of Kason 148 Anjanasakaraj). But in Thai Buddhist tradition, it was 11 March 545 BCE, the date which the current Thai lunisolar and solar calendars use as the epochal date. Yet, the Thai calendars for some reason have fixed the difference between their Buddhist Era (BE) numbering and the Christian/Common Era (CE) numbering at 543, which points to an epochal year of 544 BCE, not 545 BCE. In Myanmar, the difference between BE and CE can be 543 or 544 for CE dates, and 544 or 543 for BCE dates, depending on the month of the Buddhist Era (as the Buddhist calendar straddles the Gregorian calendar—currently from April to April).
Month
Types
The calendar recognizes two types of months: synodic month and sidereal month. The Synodic months are used to compose the years while the 27 lunar sidereal days (Sanskrit: nakshatra), alongside the 12 signs of the zodiac, are used for astrological calculations. (The Burmese calendar also recognizes a solar month called Thuriya Matha, which is defined as 1/12th of a year. But the solar month varies by the type of year such as tropical year, sidereal year, etc.)
Waxing and waning
The days of the month are counted in two halves, waxing and waning. The 15th of the waxing is the civil full moon day. The civil new moon day is the last day of the month (14th or 15th waning). Because of the inaccuracy of the calendrical calculation systems, the mean and real (true) New Moons rarely coincide. The mean New Moon often precedes the real New Moon.
Number of days per month
As the Synodic lunar month is approximately 29.5 days, the calendar uses alternating months of 29 and 30 days.
Month numbering
Various regional versions of Chula Sakarat/Burmese calendar existed across various regions of mainland Southeast Asia. Unlike Burmese systems, Kengtung, Sipsongpanna, Lan Na, Lan Xang and Sukhothai systems refer to the months by numbers, not by names. This means reading ancient texts and inscriptions in Thailand requires constant vigilance, not just in making sure one is correctly operating for the correct region, but also for variations within regions itself when incursions cause a variation in practice.
Year
The Buddhist calendar is a lunisolar calendar in which the months are based on lunar months and years are based on solar years. One of its primary objectives is to synchronize the lunar part with the solar part. The lunar months, normally twelve of them, consist alternately of 29 days and 30 days, such that a normal lunar year will contain 354 days, as opposed to the solar year of ~365.25 days. Therefore, some form of addition to the lunar year (of intercalation) is necessary. The overall basis for it is provided by cycles of 57 years. Eleven extra days are inserted in every 57 years, and seven extra months of 30 days are inserted in every 19 years (21 months in 57 years). This provides 20819 complete days to both calendars. This 57-year cycle would provide a mean year of about 365.2456 days and a mean month of about 29.530496 days, if not corrected.
As such, the calendar adds an intercalary month in leap years and sometimes also an intercalary day in great leap years. The intercalary month not only corrects the length of the year but also corrects the accumulating error of the month to extent of half a day. The average length of the month is further corrected by adding a day to Nayon at irregular intervals—a little more than seven times in two cycles (39 years). The intercalary day is never inserted except in a year which has an intercalary month. The Hindu calendar inserts an intercalary month at any time of year as soon as the accumulated fractions amount to one month. The Burmese calendar however always inserts the intercalary month at the same time of the year, after the summer solstice while the Arakanese calendar inserts it after the vernal equinox.
Burmese
The Burmese calendar year consists of 354, 384 or 385 days.
Note: The Arakanese calendar adds the intercalary day in Tagu, not in Nayon.
Cambodian, Lao and Thai
The Cambodian, Lao and Thai lunisolar calendars use a slightly different method to place the intercalary day. Instead of it in a leap year as in the Burmese system, the Thai system places it in a separate year. Thus, the Thai small leap year has 355 days while the Thai great leap year has 384 days.
New Year's Day
Since the main purpose of Buddhist calendar is to keep pace with the solar year, the new year is always marked by the solar year, which falls at the time when the Sun enters Aries. The date, which at the present falls on the 17th of April, has slowly drifted over the centuries. In the 20th century, the New Year's Day fell on April 15 or 16th but in the 17th century, it fell on April 9 or 10th. Thailand and Cambodia no longer use the traditional lunisolar calendar to mark the New Year's Day.
Cycle
The Cambodian, Lao and Thai systems give animal names to the years from a cycle of 12. The practice also existed in Burma in the Pagan period but later died out.
The Cambodian calendar also maintains a 10-year naming cycle (numbered one to ten). Cambodians use multiple systems to identify a given year. For instance, 2017 is identified as 2561 Buddhist Era, the year of Rooster, Nuppasak (Year 9). The Thai lunar calendar also uses a similar numbered 10-year cycle. Each number in the cycle corresponds to the last digit of the year in the Chula Sakarat calendar.
Accuracy
The Southeast Asian Buddhist calendars use lunar months but try to keep pace with the solar year, by inserting intercalary months and days on the Metonic cycle (in the case of the Burmese calendar, on a modified Metonic cycle). However, the solar year as defined by the Buddhist calendars is really a sidereal year, which is nearly 24 minutes longer than the actual mean tropical year. Therefore, like all sidereal-based calendars, the lunisolar calendars are slowly drifting away from the seasons. The calendars are drifting one day approximately every 60 years and 4 months.
The accumulating drift against the seasons means the New Year's Day which used to fall on 22 March (near the vernal equinox) in 638 CE now falls on 17 April in 2013 CE. There is no known internationally concerted effort to stop this drift. Thailand has moved its "Buddhist Era" to the Gregorian calendar under the name of Thai solar calendar. In Myanmar, Burmese calendarists have tried to deal with the issue by periodically modifying the intercalation schedule in the Metonic cycle. One major downside of this approach is that it is not possible to publish future calendars more than a few years (often even a year) ahead.
History
The Buddhist Era was first introduced to Southeast Asia along with Buddhism in the early centuries CE. It was not a separate calendar but simply a year numbering system that employed the organization and calculation methods of the prevailing lunisolar calendars in use throughout the region. In the early centuries CE, the reference civil calendar of the Buddhist calendar prevalent in Southeast Asia was the Saka Era (Mahāsakaraj Era), which is said to have been adopted by the Pyu state of Sri Ksetra in 80 CE. The Saka Era was gradually replaced by the Burmese Era or Culāsakaraj, first in Myanmar in 640 CE, and in other Theravada kingdoms of Southeast Asia between the 13th and 16th centuries. Theravada Buddhist tradition also recognizes pre-Buddhist Anjana Sakaraj (Añjana's Era) since the events of the Buddha's life are recorded in that era.
The tradition of using different reference calendars continued in Siam in 1912 when King Vajiravudh decreed that the Buddhist Era would now track the Thai solar calendar, the Siamese version of the Gregorian calendar with the New Year's Day of 1 April. Therefore, the Thai Buddhist Era year of 2455 began on 1 April 1912 (as opposed to 15 April 1912 according to the lunisolar calendar). The Thai Buddhist Era was further realigned to the Gregorian calendar on 6 September 1940 when Prime Minister Phibunsongkhram decreed 1 January 1941 as the start of the year 2484 BE. As a result, the Year 2483 was only 9 months long, and the Thai Buddhist Era equals that of the Common Era plus 543 years.
Current usage
The lunisolar calendar is used to mark important Buddhist holidays. Many of the holidays are celebrated as public holidays.
Computer support
The Thai-style "Buddhist calendar", which is the Gregorian calendar with the Buddhist era, is supported in Java 8, iOS, and macOS.
| Technology | Calendars | null |
1227579 | https://en.wikipedia.org/wiki/Nigerian%20Dwarf%20goat | Nigerian Dwarf goat | The Nigerian Dwarf is a modern American breed of dwarf goat. Like the American Pygmy Goat, it derives from the West African Dwarf group of breeds of West Africa.
History
Between about 1930 and 1960 a variety of small goats of the West African Dwarf group of breeds were imported from Africa to the United States to be exhibited in zoos. The Nigerian Dwarf, like the American Pygmy Goat, derives from these, but does not resemble the stocky West African Dwarf in conformation – it has been bred to have the appearance of a miniature dairy goat. It was at first reared as a show breed and companion animal; selection was for appearance and for docility. It was later found to be suitable for small-scale dairy production, and some breeding was directed towards dairy qualities. A herd-book was established in 1980.
Numbers grew rapidly; by 2002 there were almost head registered. The breed was recognized by the American Dairy Goat Association in 2005. The Nigerian Dwarf was formerly listed on the heritage breeds watchlist of the Livestock Conservancy as "recovering", but was removed from the list in 2013.
In 2023 the total number of the goats in the United States was just under ; populations numbering head are reported by Australia and Canada. In 2024 the international conservation status of the breed was "not at risk".
In the 1990s the Nigora breed was created by cross-breeding the Nigerian Dwarf with Angora and other mohair breeds.
Characteristics
The Nigerian Dwarf is small but well-proportioned; its conformation resembles that of larger dairy goats. It may be horned or naturally hornless. The coat is fine and fairly short, and may be of any color, or multicolored; common colors are gold, chocolate and black, frequently with white markings. The facial profile may be concave or straight; the ears are upright. The average weight is approximately , while maximum height is about for males and slightly less for females; stock bred for dairy performance may be rather larger than show or companion animals.
It is a precocious breed – young stock may be bred from an early age: males from about three months, females from seven or eight months. The gestation period is in the range 145–153 days; the twinning rate is high, and triplet and quadruplet births are not uncommon. Life expectancy is from eight to 12 years.
Use
The Nigerian Dwarf was originally bred for show and as a companion animal. It was later also bred for dairy use. Average milk yield of dairy stock is per year; a yield of in a lactation of 305 days was recorded in 2018. Lactation usually lasts for about ten months. The milk is high in butterfat and protein, averaging 6.5% and 3.9% respectively, and is suitable for making cheese and butter.
| Biology and health sciences | Goats | Animals |
1228141 | https://en.wikipedia.org/wiki/Dry%20season | Dry season | The dry season is a yearly period of low rainfall, especially in the tropics. The weather in the tropics is dominated by the tropical rain belt, which moves from the northern to the southern tropics and back over the course of the year. The temperate counterpart to the tropical dry season is summer or winter.
Rain belt
The tropical rain belt lies in the southern hemisphere roughly from November to March; during that time the northern tropics have a dry season with sparser precipitation, and days are typically sunny throughout. From May to September, the rain belt lies in the northern hemisphere, and the southern tropics have their dry season. Under the Köppen climate classification, for tropical climates, a dry season month is defined as a month when average precipitation is below .
The rain belt reaches roughly as far north as the Tropic of Cancer and as far south as the Tropic of Capricorn. Near these latitudes, there is one wet season and one dry season annually. At the equator there are two wet and two dry seasons, as the rain belt passes over twice a year, once moving north and once moving south. Between the tropics and the equator, locations may experience a short wet or a long wet season; and a short dry or a long dry season. Local geography may substantially modify these climate patterns, however they can.
Drought
During the dry season, humidity is very low, causing some watering holes and rivers to dry up. This lack of water (and lack of food supply) may force many grazing animals to migrate to more fertile spots. Examples of such animals are: zebras, elephants, giraffes, hippos, rhinos, antelopes and wildebeest, wild water buffaloes, African buffaloes, gaur, tapirs, emus, ostriches, rheas, and kangaroos. Because of the lack of water in the plants, bushfires (wildfires) are common.
Diseases
Data shows that in Africa the start of the dry season coincides with a rise in the cases of measles—which researchers believe might be attributed to the higher concentration of people in the dry season, as agricultural operations are all but impossible without irrigation. During this time, some farmers move into cities, creating hubs of higher population density, and allowing the disease to spread more easily.
Research
New data shows that in the seasonal parts of the South American Amazon rainforest, foliage growth and coverage varies between the dry and wet seasons—with about 25% more leaves and faster growth in the dry season. Researchers believe that the Amazon itself has an effect in bringing the onset of the wet season: by growing more foliage, it evaporates more water. However, this growth appears only in the undisturbed parts of the Amazon basin, where researchers believe roots can reach deeper and gather more rainwater. It has also been shown that ozone levels are much higher in the dry than in the wet season in the Amazon basin.
| Physical sciences | Seasons | Earth science |
1228152 | https://en.wikipedia.org/wiki/Wet%20season | Wet season | The wet season (sometimes called the rainy season or monsoon season) is the time of year when most of a region's average annual rainfall occurs. Generally, the season lasts at least one month. The term green season is also sometimes used as a euphemism by tourist authorities. Areas with wet seasons are dispersed across portions of the tropics and subtropics.
Under the Köppen climate classification, for tropical climates, a wet season month is defined as a month where average precipitation is or more. In contrast to areas with savanna climates and monsoon regimes, Mediterranean climates have wet winters and dry summers. Dry and rainy months are characteristic of tropical seasonal forests: in contrast to tropical rainforests, which do not have dry or wet seasons, since their rainfall is equally distributed throughout the year. Some areas with pronounced rainy seasons will see a break in rainfall mid-season, when the Intertropical Convergence Zone or monsoon trough moves to higher latitudes in the middle of the warm season.
When the wet season occurs during a warm season, or summer, precipitation falls mainly during the late afternoon and early evening. In the wet season, air quality improves, fresh water quality improves, and vegetation grows substantially, leading to crop yields late in the season. Rivers overflow their banks, and some animals retreat to higher ground. Soil nutrients diminish and erosion increases. The incidence of malaria and dengue increases in areas where the rainy season coincides with high temperatures, particularly in tropical areas. Some animals have adaptation and survival strategies for the wet season. Often, the previous dry season leads to food shortages in the wet season, as the crops have yet to mature. Crops which can be successfully planted during the wet or rainy season are cassava, maize, groundnut, millet, rice and yam.
The temperate counterpart to the tropical wet season is spring or autumn.
Character of the rainfall
In areas where the heavy rainfall is associated with a wind shift, the wet season is known as the monsoon season. Many tropical and subtropical climates experience monsoon rainfall patterns. Rainfall in the wet season is mainly due to daytime heating, which leads to diurnal thunderstorm activity within a pre-existing moist airmass , so the rain mainly falls in late afternoon and early evening in savanna and monsoon regions.
Much of the total rainfall each day occurs in the first minutes of the downpour, before the storms mature into their stratiform stage. Most places have only one wet season, but areas of the tropics can have two wet seasons, because the monsoon trough, or Intertropical Convergence Zone, can pass over locations in the tropics twice per year. However, since rain forests have rainfall spread evenly through the year, they do not have a wet season.
Areas affected
Areas with a savanna climate in Sub-Saharan Africa, such as Ghana, Burkina Faso, Darfur, Eritrea, Ethiopia, and Botswana have a distinct rainy season. Also subtropical areas like Florida, South and Southeast Texas, and southern Louisiana in the United States have a rainy season. Monsoon regions include the Indian subcontinent, Southeast Asia (including Indonesia and Philippines), northern sections of Australia, Polynesia, Central America, western and southern Mexico, the Desert Southwest of the United States, southern Guyana, and northeast Brazil.
Northern Guyana has two wet seasons: one in early spring and the other in early winter. In western Africa, there are two rainy seasons across southern sections, but only one across the north. Within the Mediterranean climate regime, the west coast of the United States, the southwest coast of Australia and South Africa, the Mediterranean coastline of Italy, Spain, Greece, Lebanon, Syria, Algeria, Morocco, Tunisia, and Turkey, as well as areas further inland in Western Asia which include Jordan, Northern Iraq and most parts of Iran, experience a wet season in the winter months. Similarly, the wet season in the Negev Desert of Israel extends from October through May. At the boundary between the Mediterranean and monsoon climates lies the Sonoran Desert, which receives the two rainy seasons associated with each climate regime.
The wet season is known by many different local names throughout the world. For example, in Mexico it is known as "storm season". Different names are given to the various short "seasons" of the year by the First Nations of Northern Australia: the wet season typically experienced there from December to March is called Gudjewg. The precise meaning of the word is disputed, although it is widely accepted to relate to the severe thunderstorms, flooding, and abundant vegetation growth commonly experienced at this time.
Effects
In tropical areas, when the monsoon arrives, high daytime high temperatures drop and overnight low temperatures increase, thus reducing diurnal temperature variation. During the wet season, a combination of heavy rainfall and, in some places such as Hong Kong, an onshore wind, improve air quality.
In Brazil, the wet season is correlated with weaker trade winds off the ocean. The pH level of water becomes more balanced due to the charging of local aquifers during the wet season. Water also softens, as the concentration of dissolved materials reduces during the rainy season. Erosion is also increased during rainy periods.
Arroyos that are dry at other times of the year fill with runoff, in some cases with water as deep as . Leaching of soils during periods of heavy rainfall depletes nutrients. The higher runoff from land masses affects nearby ocean areas, which are more stratified, or less mixed, due to stronger surface currents forced by the heavy rainfall runoff.
Floods
High rainfall can cause widespread flooding, which can lead to landslides and mudflows in mountainous areas. Such floods cause rivers to burst their banks and submerge homes. The Ghaggar-Hakra River, which only flows during India's monsoon season, can flood and severely damage local crops. Floods can be exacerbated by fires that occurred during the previous dry season, which cause soils which are sandy or composed of loam to become hydrophobic, or water repellent. In various ways governments may help people deal with wet season floods. Flood plain mapping identifies which areas are more prone to flooding. Instructions on controlling erosion through outreach are also provided by telephone or the internet.
Life adaptations
Humans
The wet season is the main period of vegetation growth within the Savanna climate regime. However, this also means that wet season is a time for food shortages before crops reach their full maturity. This causes seasonal weight changes for people in developing countries, with a drop occurring during the wet season until the time of the first harvest, when weights rebound. Malaria incidence increases during periods of high temperature and heavy rainfall.
Animals
Cows calve, or give birth, at the beginning of the wet season. The onset of the rainy season signals the departure of the monarch butterfly from Mexico. Tropical species of butterflies show larger dot markings on their wings to fend off possible predators and are more active during the wet season than the dry season. Within the tropics and warmer areas of the subtropics, decreased salinity of near shore wetlands due to the rains causes an increase in crocodile nesting. Other species, such as the arroyo toad, spawn within the couple of months after the seasonal rains. Armadillos and rattlesnakes seek higher ground.
| Physical sciences | Seasons | Earth science |
1228320 | https://en.wikipedia.org/wiki/Energy%20return%20on%20investment | Energy return on investment | In energy economics and ecological energetics, energy return on investment (EROI), also sometimes called energy returned on energy invested (ERoEI), is the ratio of the amount of usable energy (the exergy) delivered from a particular energy resource to the amount of exergy used to obtain that energy resource.
Arithmetically the EROI can be defined as:
.
When the EROI of a source of energy is less than or equal to one, that energy source becomes a net "energy sink", and can no longer be used as a source of energy. A related measure, called energy stored on energy invested (ESOEI), is used to analyse storage systems.
To be considered viable as a prominent fuel or energy source a fuel or energy must have an EROI ratio of at least 3:1.
History
The energy analysis field of study is credited with being popularized by Charles A. S. Hall, a Systems ecology and biophysical economics professor at the State University of New York. Hall applied the biological methodology, developed at an Ecosystems Marine Biological Laboratory, and then adapted that method to research human industrial civilization. The concept would have its greatest exposure in 1984, with a paper by Hall that appeared on the cover of the journal Science.
Application to various technologies
Photovoltaic
The issue is still subject of numerous studies, and prompting academic argument. That's mainly because the "energy invested" critically depends on technology, methodology, and system boundary assumptions, resulting in a range from a maximum of 2000 kWh/m2 of module area down to a minimum of 300 kWh/m2 with a median value of 585 kWh/m2 according to a meta-study from 2013.
Regarding output, it obviously depends on the local insolation, not just the system itself, so assumptions have to be made.
Some studies (see below) include in their analysis that photovoltaic produce electricity, while the invested energy may be lower grade primary energy.
A 2015 review in Renewable and Sustainable Energy Reviews assessed the energy payback time and EROI of a variety of PV module technologies. In this study, which uses an insolation of 1700 kWh/m2/yr and a system lifetime of 30 years, mean harmonized EROIs between 8.7 and 34.2 were found. Mean harmonized energy payback time varied from 1.0 to 4.1 years.
In 2021, the Fraunhofer Institute for Solar Energy Systems calculated an energy payback time of around 1 year for European PV installations (0.9 years for Catania in Southern Italy, 1.1 years for Brussels) with wafer-based silicon PERC cells.
Wind turbines
In the scientific literature EROIs wind turbines is around 16 unbuffered and 4 buffered. Data collected in 2018 found that the EROI of operational wind turbines averaged 19.8 with high variability depending on wind conditions and wind turbine size.
EROIs tend to be higher for recent wind turbines compared to older technology wind turbines.
Vestas reports an EROI of 31 for its V150 model wind turbine.
Hydropower plants
The EROI for hydropower plants averages to about 110 when it is run for about 100 years.
Oil sands
Because much of the energy required for producing oil from oil sands (bitumen) comes from low value fractions separated out by the upgrading process, there are two ways to calculate EROI, the higher value given by considering only the external energy inputs and the lower by considering all energy inputs, including self generated. One study found that in 1970 oil sands net energy returns was about 1.0 but by 2010 had increased to about 5.23.
Conventional oil
Conventional sources of oil have a rather large variation depending on various geologic factors. The EROI for refined fuel from conventional oil sources varies from around 18 to 43.
Oil Shale
Due to the process heat input requirements for oil shale harvesting, the EROI is low. Typically natural gas is used, either directly combusted for process heat or used to power an electricity generating turbine, which then uses electrical heating elements to heat the underground layers of shale to produce oil from the kerogen. Resulting EROI is typically around 1.4-1.5. Economically, oil shale might be viable due to the effectively free natural gas on site used for heating the kerogen, but opponents have debated that the natural gas could be extracted directly and used for relatively inexpensive transportation fuel rather than heating shale for a lower EROI and higher carbon emissions.
Oil liquids
The weighted average standard EROI of all oil liquids (including coal-to-liquids, gas-to-liquids, biofuels, etc.) is expected to decrease from 44.4 in 1950 to a plateau of 6.7 in 2050.
Natural gas
The standard EROI for natural gas is estimated to decrease from 141.5 in 1950 to an apparent plateau of 16.8 in 2050.
Nuclear plants
The EROI for nuclear plants ranges from 20 to 81.
Non-manmade energy inputs
The natural or primary energy sources are not included in the calculation of energy invested, only the human-applied sources.
For example, in the case of biofuels the solar insolation driving photosynthesis is not included, and the energy used in the stellar synthesis of fissile elements is not included for nuclear fission. The energy returned includes only human usable energy and not wastes such as waste heat.
Nevertheless, heat of any form can be counted where it is actually used for heating. However the use of waste heat in district heating and water desalination in cogeneration plants is rare, and in practice it is often excluded in EROI analysis of energy sources.
Competing methodology
In a 2010 paper by Murphy and Hall, the advised extended ["Ext"] boundary protocol, for all future research on EROI, was detailed. In order to produce, what they consider, a more realistic assessment and generate greater consistency in comparisons, than what Hall and others view as the "weak points" in a competing methodology. In more recent years, however, a source of continued controversy is the creation of a different methodology endorsed by certain members of the IEA which for example most notably in the case of photovoltaic solar panels, controversially generates more favorable values.
In the case of photovoltaic solar panels, the IEA method tends to focus on the energy used in the factory process alone. In 2016, Hall observed that much of the published work in this field is produced by advocates or persons with a connection to business interests among the competing technologies, and that government agencies had not yet provided adequate funding for rigorous analysis by more neutral observers.
Relationship to net energy gain
EROI and Net energy (gain) measure the same quality of an energy source or sink in numerically different ways. Net energy describes the amounts, while EROI measures the ratio or efficiency of the process. They are related simply by
or
For example, given a process with an EROI of 5, expending 1 unit of energy yields a net energy gain of 4 units. The break-even point happens with an EROI of 1 or a net energy gain of 0. The time to reach this break-even point is called energy payback period (EPP) or energy payback time (EPBT).
Economic influence
Although many qualities of an energy source matter (for example oil is energy-dense and transportable, while wind is variable), when the EROI of the main sources of energy for an economy fall that energy becomes more difficult to obtain and its relative price may increase.
In regard to fossil fuels, when oil was originally discovered, it took on average one barrel of oil to find, extract, and process about 100 barrels of oil. The ratio, for discovery of fossil fuels in the United States, has declined steadily over the last century from about 1000:1 in 1919 to only 5:1 in the 2010s.
Since the invention of agriculture, humans have increasingly used exogenous sources of energy to multiply human muscle-power.
Some historians have attributed this largely to more easily exploited (i.e. higher EROI) energy sources, which is related to the concept of energy slaves. Thomas Homer-Dixon argues that a falling EROI in the Later Roman Empire was one of the reasons for the collapse of the Western Empire in the fifth century CE. In "The Upside of Down" he suggests that EROI analysis provides a basis for the analysis of the rise and fall of civilisations. Looking at the maximum extent of the Roman Empire, (60 million) and its technological base the agrarian base of Rome was about 1:12 per hectare for wheat and 1:27 for alfalfa (giving a 1:2.7 production for oxen). One can then use this to calculate the population of the Roman Empire required at its height, on the basis of about 2,500–3,000 calories per day per person. It comes out roughly equal to the area of food production at its height. But ecological damage (deforestation, soil fertility loss particularly in southern Spain, southern Italy, Sicily and especially north Africa) saw a collapse in the system beginning in the 2nd century, as EROI began to fall. It bottomed in 1084 when Rome's population, which had peaked under Trajan at 1.5 million, was only 15,000.
Evidence also fits the cycle of Mayan and Cambodian collapse too. Joseph Tainter suggests that diminishing returns of the EROI is a chief cause of the collapse of complex societies, which has been suggested as caused by peak wood in early societies. Falling EROI due to depletion of high quality fossil fuel resources also poses a difficult challenge for industrial economies, and could potentially lead to declining economic output and challenge the concept (which is very recent when considered from a historical perspective) of perpetual economic growth.
Criticism of EROI
EROI is calculated by dividing the energy output by the energy input. Measuring total energy output is often easy, especially in the case for an electrical output where some appropriate electricity meter can be used. However, researchers disagree on how to determine energy input accurately and therefore arrive at different numbers for the same source of energy.
How deep should the probing in the supply chain of the tools being used to generate energy go? For example, if steel is being used to drill for oil or construct a nuclear power plant, should the energy input of the steel be taken into account? Should the energy input into building the factory being used to construct the steel be taken into account and amortized? Should the energy input of the roads which are used to ferry the goods be taken into account? What about the energy used to cook the steelworkers' breakfasts? These are complex questions evading simple answers. A full accounting would require considerations of opportunity costs and comparing total energy expenditures in the presence and absence of this economic activity.
However, when comparing two energy sources a standard practice for the supply chain energy input can be adopted. For example, consider the steel, but don't consider the energy invested in factories deeper than the first level in the supply chain. It is in part for these fully encompassed systems reasons, that in the conclusions of Murphy and Hall's paper in 2010, an EROI of 5 by their extended methodology is considered necessary to reach the minimum threshold of sustainability, while a value of 12–13 by Hall's methodology is considered the minimum value necessary for technological progress and a society supporting high art.
Richards and Watt propose an Energy Yield Ratio for photovoltaic systems as an alternative to EROI (which they refer to as Energy Return Factor). The difference is that it uses the design lifetime of the system, which is known in advance, rather than the actual lifetime. This also means that it can be adapted to multi-component systems where the components have different lifetimes.
Another issue with EROI that many studies attempt to tackle is that the energy returned can be in different forms, and these forms can have different utility. For example, electricity can be converted more efficiently than thermal energy into motion, due to electricity's lower entropy. In addition, the form of energy of the input can be completely different from the output. For example, energy in the form of coal could be used in the production of ethanol. This might have an EROI of less than one, but could still be desirable due to the benefits of liquid fuels (assuming the latters are not used in the processes of extraction and transformation).
Additional EROI calculations
There are three prominent expanded EROI calculations, they are point of use, extended and societal. Point of Use EROI expands the calculation to include the cost of refining and transporting the fuel during the refining process. Since this expands the bounds of the calculation to include more production process EROI will decrease. Extended EROI includes point of use expansions as well as including the cost of creating the infrastructure needed for transportation of the energy or fuel once refined. Societal EROI is a sum of all the EROIs of all the fuels used in a society or nation. A societal EROI has never been calculated and researchers believe it may currently be impossible to know all variables necessary to complete the calculation, but attempted estimates have been made for some nations. Calculations are done by summing all of the EROIs for domestically produced and imported fuels and comparing the result to the Human Development Index (HDI), a tool often used to understand well-being in a society. According to this calculation, the amount of energy a society has available to them increases the quality of life for the people living in that country, and countries with less energy available also have a harder time satisfying citizens' basic needs. This is to say that societal EROI and overall quality of life are very closely linked.
EROI and payback periods of some types of power plants
The following table is a compilation of sources of energy. The minimum requirement is a breakdown of the cumulative energy expenses according to material data. Frequently in literature harvest factors are reported, for which the origin of the values is not completely transparent. These are not included in this table.
The bold numbers are those given in the respective literature source, the normal printed ones are derived (see Mathematical Description).
(a) The cost of fuel transportation is taken into account
(b) The values refer to the total energy output. The expense for storage power plants, seasonal reserves or conventional load balancing power plants is not taken into account.
(c) The data for the E-82 come from the manufacturer, but are confirmed by TÜV Rheinland.
ESOEI
ESOEI (or ESOIe) is used when EROI is below 1. "ESOIe is the ratio of electrical energy stored over the lifetime of a storage device to the amount of embodied electrical energy required to build the device."
One of the notable outcomes of the Stanford University team's assessment on ESOI, was that if pumped storage was not available, the combination of wind energy and the commonly suggested pairing with battery technology as it presently exists, would not be sufficiently worth the investment, suggesting instead curtailment.
EROI under rapid growth
A related recent concern is energy cannibalism where energy technologies can have a limited growth rate if climate neutrality is demanded. Many energy technologies are capable of replacing significant volumes of fossil fuels and concomitant green house gas emissions. Unfortunately, neither the enormous scale of the current fossil fuel energy system nor the necessary growth rate of these technologies is well understood within the limits imposed by the net energy produced for a growing industry. This technical limitation is known as energy cannibalism and refers to an effect where rapid growth of an entire energy producing or energy efficiency industry creates a need for energy that uses (or cannibalizes) the energy of existing power plants or production plants.
The overcomes some of these problems. A solar breeder is a photovoltaic panel manufacturing plant which can be made energy-independent by using energy derived from its own roof using its own panels. Such a plant becomes not only energy self-sufficient but a major supplier of new energy, hence the name solar breeder. Research on the concept was conducted by Centre for Photovoltaic Engineering, University of New South Wales, Australia. The reported investigation establishes certain mathematical relationships for the solar breeder which clearly indicate that a vast amount of net energy is available from such a plant for the indefinite future. The solar module processing plant at Frederick, Maryland was originally planned as such a solar breeder. In 2009 the Sahara Solar Breeder Project was proposed by the Science Council of Japan as a cooperation between Japan and Algeria with the highly ambitious goal of creating hundreds of GW of capacity within 30 years.
| Technology | Concepts | null |
1228679 | https://en.wikipedia.org/wiki/Gyromagnetic%20ratio | Gyromagnetic ratio | In physics, the gyromagnetic ratio (also sometimes known as the magnetogyric ratio in other disciplines) of a particle or system is the ratio of its magnetic moment to its angular momentum, and it is often denoted by the symbol , gamma. Its SI unit is the radian per second per tesla (rad⋅s−1⋅T−1) or, equivalently, the coulomb per kilogram (C⋅kg−1).
The term "gyromagnetic ratio" is often used as a synonym for a different but closely related quantity, the -factor. The -factor only differs from the gyromagnetic ratio in being dimensionless.
For a classical rotating body
Consider a nonconductive charged body rotating about an axis of symmetry. According to the laws of classical physics, it has both a magnetic dipole moment due to the movement of charge and an angular momentum due to the movement of mass arising from its rotation. It can be shown that as long as its charge and mass density and flow are distributed identically and rotationally symmetric, its gyromagnetic ratio is
where is its charge and is its mass.
The derivation of this relation is as follows. It suffices to demonstrate this for an infinitesimally narrow circular ring within the body, as the general result then follows from an integration. Suppose the ring has radius , area , mass , charge , and angular momentum . Then the magnitude of the magnetic dipole moment is
For an isolated electron
An isolated electron has an angular momentum and a magnetic moment resulting from its spin. While an electron's spin is sometimes visualized as a literal rotation about an axis, it cannot be attributed to mass distributed identically to the charge. The above classical relation does not hold, giving the wrong result by the absolute value of the electron's -factor, which is denoted :
where is the Bohr magneton.
The gyromagnetic ratio due to electron spin is twice that due to the orbiting of an electron.
In the framework of relativistic quantum mechanics,
where is the fine-structure constant. Here the small corrections to the relativistic result come from the quantum field theory calculations of the anomalous magnetic dipole moment. The electron -factor is known to twelve decimal places by measuring the electron magnetic moment in a one-electron cyclotron:
The electron gyromagnetic ratio is
The electron -factor and are in excellent agreement with theory; see Precision tests of QED for details.
Gyromagnetic factor not as a consequence of relativity
Since a gyromagnetic factor equal to 2 follows from Dirac's equation, it is a frequent misconception to think that a -factor 2 is a consequence of relativity; it is not. The factor 2 can be obtained from the linearization of both the Schrödinger equation and the relativistic Klein–Gordon equation (which leads to Dirac's). In both cases a 4-spinor is obtained and for both linearizations the -factor is found to be equal to 2. Therefore, the factor 2 is a consequence of the minimal coupling and of the fact of having the same order of derivatives for space and time.
Physical spin- particles which cannot be described by the linear gauged Dirac equation satisfy the gauged Klein–Gordon equation extended by the term according to,
Here, and stand for the Lorentz group generators in the Dirac space, and the electromagnetic tensor respectively, while is the electromagnetic four-potential. An example for such a particle is the spin companion to spin in the representation space of the Lorentz group. This particle has been shown to be characterized by and consequently to behave as a truly quadratic fermion.
For a nucleus
Protons, neutrons, and many nuclei carry nuclear spin, which gives rise to a gyromagnetic ratio as above. The ratio is conventionally written in terms of the proton mass and charge, even for neutrons and for other nuclei, for the sake of simplicity and consistency. The formula is:
where is the nuclear magneton, and is the -factor of the nucleon or nucleus in question. The ratio equal to , is 7.622593285(47) MHz/T.
The gyromagnetic ratio of a nucleus plays a role in nuclear magnetic resonance (NMR) and magnetic resonance imaging (MRI). These procedures rely on the fact that bulk magnetization due to nuclear spins precess in a magnetic field at a rate called the Larmor frequency, which is simply the product of the gyromagnetic ratio with the magnetic field strength. With this phenomenon, the sign of determines the sense (clockwise vs counterclockwise) of precession.
Most common nuclei such as 1H and 13C have positive gyromagnetic ratios. Approximate values for some common nuclei are given in the table below.
Larmor precession
Any free system with a constant gyromagnetic ratio, such as a rigid system of charges, a nucleus, or an electron, when placed in an external magnetic field (measured in teslas) that is not aligned with its magnetic moment, will precess at a frequency (measured in hertz) proportional to the external field:
For this reason, values of , in units of hertz per tesla (Hz/T), are often quoted instead of .
Heuristic derivation
The derivation of this ratio is as follows: First we must prove the torque resulting from subjecting a magnetic moment to a magnetic field is
The identity of the functional form of the stationary electric and magnetic fields has led to defining the magnitude of the magnetic dipole moment equally well as , or in the following way, imitating the moment of an electric dipole: The magnetic dipole can be represented by a needle of a compass with fictitious magnetic charges on the two poles and vector distance between the poles under the influence of the magnetic field of earth By classical mechanics the torque on this needle is But as previously stated so the desired formula comes up. is the unit distance vector.
The spinning electron model here is analogous to a gyroscope. For any rotating body the rate of change of the angular momentum equals the applied torque :
Note as an example the precession of a gyroscope. The earth's gravitational attraction applies a force or torque to the gyroscope in the vertical direction, and the angular momentum vector along the axis of the gyroscope rotates slowly about a vertical line through the pivot. In place of a gyroscope, imagine a sphere spinning around the axis with its center on the pivot of the gyroscope, and along the axis of the gyroscope two oppositely directed vectors both originated in the center of the sphere, upwards and downwards Replace the gravity with a magnetic flux density
represents the linear velocity of the pike of the arrow along a circle whose radius is where is the angle between and the vertical. Hence the angular velocity of the rotation of the spin is
Consequently,
This relationship also explains an apparent contradiction between the two equivalent terms, gyromagnetic ratio versus magnetogyric ratio: whereas it is a ratio of a magnetic property (i.e. dipole moment) to a gyric (rotational, from , "turn") property (i.e. angular momentum), it is also, at the same time, a ratio between the angular precession frequency (another gyric property) and the magnetic field.
The angular precession frequency has an important physical meaning: It is the angular cyclotron frequency, the resonance frequency of an ionized plasma being under the influence of a static finite magnetic field, when we superimpose a high frequency electromagnetic field.
| Physical sciences | Quantum mechanics | Physics |
1228845 | https://en.wikipedia.org/wiki/Moisturizer | Moisturizer | A moisturizer, or emollient, is a cosmetic preparation used for protecting, moisturizing, and lubricating the skin. These functions are normally performed by sebum produced by healthy skin. The word "emollient" is derived from the Latin verb mollire, to soften.
Mechanism of action
In the human body, water constantly evaporates from the deeper layers of the skin through an effect known as transepidermal water loss. By regulating its water content, human skin naturally maintains a dry, easily shed surface as a barrier against pathogens, dirt, or damage, while protecting itself from drying out and becoming brittle and rigid. The ability to retain moisture depends on the lipid bilayer between the dead skin cells.
Moisturizers modify the rate of water loss, with active ingredients of moisturizers falling into one of two categories: occlusives and humectants.
Occlusives form a hydrophobic coating on the surface of the skin, keeping moisture from escaping. The more occlusive the formulation, the greater the effect. Ointments are more occlusive than aqueous creams, which are more occlusive than lotion. Water loss through the skin is normally about 4–8 g/(m2⋅h). A layer of petrolatum applied to normal skin can reduce that loss by 50–75% for several hours. Oils naturally produced by the human body moisturize through this same mechanism.
Humectants are hydrophilic and absorb water. They absorb water from humid air (when >70% humidity) to moisturize the skin. More commonly, however, they draw out water from the dermis into the epidermis, making deeper skin dryer. When used in practical applications, humectants are combined with occlusives. Moisturizers commonly contain water, which acts as a temporary hydration agent.
Kinds
There are many different types of moisturizers. Petrolatum is one of the most effective moisturizers, although it can be unpopular due to its oily consistency.
Other popular moisturizers are cetyl alcohol, cetearyl alcohol, cocoa butter, isopropyl myristate, isopropyl palmitate, lanolin, liquid paraffin, polyethylene glycols, shea butter, silicone oils, stearic acid, stearyl alcohol and castor oil, and other oils.
Moisturizers may also be available as lotions, creams, ointments, bath oils, or soap substitutes.
Mineral oils and waxes are not prone to oxidation or rancidity. For this reason, they have essentially replaced vegetable oils in emollients and topical medication.
Moisturizer cosmetics may additionally contain antioxidants, ceramides, emulsifiers, fragrances, penetration enhancers, preservatives, and solvents. Some products are marketed as having anti-wrinkle and skin enhancement effects. Many plant and animal extracts have been claimed to impart skin benefits, but such claims are presented with little scientific evidence.
Use
Moisturizers are used for the treatment of certain skin diseases, such as psoriasis, ichthyosis vulgaris, xerosis, and pruritus in atopic dermatitis. More often, they are bases or vehicles for topical medication, such as in Whitfield's ointment. They are often combined with humectants, such as salicylic acid and urea.
Moisturizers are also widely used in sunscreens, antiperspirants, skin cleansers, shaving creams, aftershaves, and hair tonics.
Moisturizers are used in disposable diapers to prevent dry skin and diaper dermatitis.
Moisturizers show some beneficial effects in treating atopic dermatitis (eczema). Using moisturizers helps to improve skin comfort and may reduce disease flares. They can be used as leave-on treatments, bath additives, or soap substitutes. There are many different moisturizer products, but the majority of leave-on treatments (from least to most greasy) are one of the following: lotions, creams, gels, or ointments. As none of the different types of moisturizers are more effective than the others, people with atopic dermatitis need to choose one or more products according to their age, affected body site, climate/season, and personal preference. However using moisturizers daily in infants during the first year of life does not help to prevent the development of atopic dermatitis, and might even increase the risk of skin infections.
Potential health risks
Over-moisturization
Persistent moisturization to the skin via water contact may contribute to an allergic reaction or contact dermatitis. This could allow foreign objects to penetrate the skin.
Changes in the skin's normal ecological environment–either atop or within the skin–can also allow for the overgrowth of pathogens.
Allergens
Aromas or food additives in moisturizers may trigger an immune reaction, including development of an allergy.
There is currently no regulation over use of the term "hypoallergenic". In fact, some pediatric skin products marketed as hypoallergenic contained allergens.
Those with eczema are especially vulnerable to allergic reaction with lotions and creams, as their compromised skin barrier allows preservatives to bind with and activate immune cells.
The American Academy of Allergy, Asthma, and Immunology released a warning in 2014 that natural lotion containing ingredients commonly found in food (such as goats milk, cow's milk, coconut milk, or oil) may introduce new allergies, potentially causing an allergic reaction upon later consuming such foods. A paper published in 2021 noted that frequent skin moisturization in early life might promote the development of food allergy, even when skin conditions such as eczema are taken into account.
Fire risk
Paraffin-based skincare products and contaminated clothing can pose a serious fire hazard. Between 2010 and 2018, paraffin was linked to 50 fire incidents (49 of which were fatal) in the U.K. A West Yorkshire Fire and Rescue Service study found that clothing contaminated with cream containing only 21% paraffin, when set alight, was fully engulfed in flame in 3 seconds. The Medicines and Healthcare products Regulatory Agency (MHRA) released a warning in 2008 about the flammability of paraffin-based products. MHRA recommends that regular paraffin users change their sheets regularly, refrain from smoking and/or bringing open flames near paraffin-coated people or objects. MHRA also recommends that skin creams containing any paraffin have a flammability warning on the packaging.
Brands of moisturizers
Artistry
Aveeno
Bath & Body Works
CeraVe
Cetaphil
Curél
Dial
Diprobase
Eucerin
Gold Bond
Garnier
Jergens
Johnson & Johnson
Lux
Neutrogena
Nivea
Olay
Sebamed
Simple Skincare
Suave
Vaseline
| Biology and health sciences | Hygiene products | Health |
1229942 | https://en.wikipedia.org/wiki/Port%20of%20Singapore | Port of Singapore | The Port of Singapore is the collection of facilities and terminals that conduct maritime trade and handle Singapore's harbours and shipping. It has been ranked as the top maritime capital of the world, since 2015. Currently the world's second-busiest port in terms of total shipping tonnage, it also transships a fifth of the world's shipping containers, half of the world's annual supply of crude oil, and is the world's busiest transshipment port. It was also the busiest port in terms of total cargo tonnage handled until 2010, when it was surpassed by the Port of Shanghai.
Because of its strategic location, Singapore has been a significant entrepôt and trading post for at least two centuries. During the contemporary era, its ports have not become just a mere economic boon for the country, but an economic necessity because Singapore is lacking in land and natural resources. The port is critical for importing natural resources, and then later re-exporting products after they have been domestically refined and shaped in some manner, for example wafer fabrication or oil refining to generate value added revenue. The Port of Singapore is also the world's largest bunkering port. The majority of ships that pass between the Indian Ocean and the Pacific Ocean go through the Singapore Strait. The Straits of Johor on the country's north are impassable for ships due to the Johor-Singapore Causeway, built in 1923, which links the town of Woodlands, Singapore to the city of Johor Bahru in Malaysia.
History
Before 1819
In the late 13th century, a Kingdom known as Singapura was established on the north bank of the Singapore River around what was called the Old Harbour. It was the only port in the southern part of the Strait of Malacca and serviced ships and traders in the region, competing with other ports along the coast of the Malacca Strait such as Jambi, Kota Cina, Lambri, Semudra, Palembang, South Kedah and Tamiang. The port had two functions. First, it made available products that were in demand by international markets; according to the Daoyi Zhilüe (Brief Annals of Foreign Islands, 1349) by Chinese trader Wang Dayuan (born 1311, fl. 1328–1339), these included top-quality hornbill casques, lakawood and cotton. Although these goods were also available from other Southeast Asian ports, those from Singapore were unique in terms of their quality. Secondly, Singapore acted as a gateway into the regional and international economic system for its immediate region. South Johor and the Riau Archipelago supplied products to Singapore for export elsewhere, while Singapore was the main source of foreign products to the region. Archaeological artefacts such as ceramics and glassware found in the Riau Archipelago evidence this. In addition, cotton was transshipped from Java or India through Singapore.
In 1984, an archaeological excavation had commenced at Fort Canning Hill led by the archaeologist Dr John Miksic. A range of artefacts including earthenware, ceramic, and porcelain pieces were found which suggests Singapore's role as an active trading port in the 14th century.
By the 15th century, Singapore had declined as an international trading port due to the ascendance of the Malacca Sultanate. Local trade continued on the island. A map of Singapore by Portuguese mathematician Manuel Godinho de Eredia showed the location of Xabandaria or the office of a shahbandar, the Malay official responsible for international trade. Shards of 15th-century Siam ceramics and late 16th – or early 17th-century Chinese blue and white porcelain have been found at the Singapore and Kallang Rivers. Singapore also provided other regional ports with local products demanded by international markets. For instance, blackwood (a generic term used by Europeans to refer to rosewood) was exported from Singapore to Malacca, and was in turn purchased by Chinese traders and shipped to China for furniture-making:)
In the early 17th century, Singapore's main settlement and its port were destroyed by a punitive force from Aceh. After this incident, there was no significant settlement or port at Singapore until 1819.
1819–1960
In 1819, Stamford Raffles, a British colonial official, excited by the deep and sheltered waters in Keppel Harbour, established for the British Empire a new settlement and international trading port on the island. Keen to attract Asian and European traders to the new port, Raffles directed that land along the banks of the Singapore River, particularly the south bank, be reclaimed where necessary and allocated to Chinese and English country traders to encourage them to establish a stake in the port-settlement. Chinese traders, because of their frequent commercial interactions with Southeast Asian traders throughout the year, set up their trading houses along the lower reaches of the river, while English country traders, who depended on the annual arrival of trade from India, set up warehouses along the upper reaches. The port relied on three main networks of trade that existed in Southeast Asia at that time: the Chinese network, which linked Southeast Asia with the southern Chinese ports of Fujian and Guangdong; the Southeast Asian network, which linked the islands of the Indonesian archipelago; and the European and Indian Ocean network, which linked Singapore to the markets of Europe and the Indian Ocean littoral. These networks were complementary, and positioned Singapore as the transshipment point of regional and international trade. By the 1830s, Singapore had overtaken Batavia (now Jakarta) as the centre of the Chinese junk trade, and also become the centre of English country trade, in Southeast Asia. This was because Southeast Asian traders preferred the free port of Singapore to other major regional ports which had cumbersome restrictions. Singapore had also supplanted Tanjung Pinang as the export gateway for the gambier and pepper industry of the Riau–Lingga Archipelago by the 1830s, and South Johor by the 1840s. It had also become the centre of the Teochew trade in marine produce and rice.
As the volume of its maritime trade increased in the 19th century, Singapore became a key port of call for sailing and steam vessels in their passage along Asian sea routes. From the 1840s, Singapore became an important coaling station for steam shipping networks that were beginning to form. Towards the late 19th century, Singapore became a staple port servicing the geographical hinterland of the Malay Peninsula. Following the institution of the British Forward Movement, Singapore became the administrative capital of British Malaya. Roads and railways were developed to transport primary materials such as crude oil, rubber and tin from the Malay Peninsula to Singapore to be processed into staple products, and then shipped to Britain and other international markets. During the colonial period, this was the most important role of the port of Singapore.
1963–2021
Singapore ceased to be part of the British Empire when it merged with Malaysia in 1963. Singapore lost its hinterland and was no longer the administrative or economic capital of the Malay Peninsula. The processing in Singapore of raw materials extracted in the Peninsula was drastically reduced due to the absence of a common market between Singapore and the Peninsular states.
Since Singapore's full independence in 1965, it has had to compete with other ports in the region to attract shipping and trade at its port. It has done so by developing an export-oriented economy based on value-added manufacturing. It obtains raw or partially manufactured products from regional and global markets and exports value-added products back to these markets through market access agreements such as World Trade Organization directives and free trade agreements.
By the 1980s, maritime trading activity had ceased in the vicinity of the Singapore River except in the form of passenger transport, as other terminals and harbours took over this role. Keppel Harbour is now home to three container terminals. Other terminals were built in Jurong and Pasir Panjang as well as in Sembawang in the north. Today, the port operations in Singapore are handled by two players: PSA International (formerly the Port of Singapore Authority) and Jurong Port, which collectively operate six container terminals and three general-purpose terminals around Singapore.
In the 1990s the port became more well-known and overtook Yokohama, and eventually became the busiest port in terms of shipping tonnage.
Singapore is part of the Maritime Silk Road that runs from the Chinese coast to the southern tip of India, to Mombasa, from there through the Red Sea via the Suez Canal to the Mediterranean, to the Upper Adriatic region of the northern Italian hub Trieste with its rail connections to Central Europe and the North Sea.
Since 2022
The Tuas Mega Port is projected to be the only port in Singapore after the PSA city terminals and Pasir Panjang Terminal are closed in 2027 and 2040 respectively, ending an era of port operations in the city area which began in 1819. The Sea Transport Industry Transformation Map (ITM) launched by the Maritime and Port Authority of Singapore (MPA) aims to grow the industry's value-add by $4.5 billion and create more than 5,000 new jobs by 2025.
Automation will be a key part of the new port, with over 1,000 battery-powered driverless vehicles and a fleet of almost 1,000 automated yard cranes to be developed for the port. Nelson Quek, PSA Singapore's head of Tuas planning stated that "Tuas, when it's fully developed, is going to be the single largest fully-automated terminal in the world". It will also be able to cater to the demands of the world's largest container ships, with 26 km of deep-water berths. Besides just handling containers, the port will have space set aside for companies to be located, a move that aims to improve the links between port and businesses. It is projected to be twice the size of Ang Mo Kio new town.
Operations at Tuas Mega Port began in September 2021, and the port officially opened on September 1, 2022 with three berths in service.
Operations
The port is the world's busiest port in terms of shipping tonnage handled, with 1.15 billion gross tons (GT) handled in 2005. In terms of cargo tonnage, Singapore is behind Shanghai with 423 million freight tons handled. The port retains its position as the world's busiest hub for transshipment traffic in 2005, and is also the world's biggest bunkering hub, with 25 million tonnes sold in the same year.
Singapore is ranked first globally in 2005 in terms of containerised traffic, with 23.2 million Twenty-foot equivalent units (TEUs) handled. High growth in containerised traffic has seen the port overtaking Hong Kong since the first quarter of 2005, and has led the race ever since, with an estimated 19,335 kTEUs handled in the year up to October, compared to 18,640 kTEUs handled in Hong Kong in the same period. A rise in regional traffic consolidating the port's position in Southeast Asia, and increases in transshipment traffic using the strategic East Asia-Europe route via Singapore helped the port to emerge tops at the end of the year, a title it had not held since overtaking Hong Kong once in 1998.
Operators
PSA Singapore's container facilities are as follows:
Container berths: 52
Quay length: 15,500 m
Area: 600 hectares
Max draft: 16 m
Quay cranes: 190
Designed capacity: 35,000 kTEU
PSA Singapore has 13 berths which are part of the Pasir Panjang Container Terminal's Phase Two which are due for completion by 2009. Phase Three and Four will add another 16 berths and are expected to be completed by 2013.
Jurong Port's facilities are as follows:
Berths: 32
Berth length: 5.6 km
Maximum vessel draft: 15.7 m
Maximum vessel size:
Area: 127 Hectares Free Trade Zone, 28 Hectares non-Free Trade Zone
Warehouse facilities: 178,000 m2
PSA Singapore also has a 40-year contract to operate the tax-free Gwadar Port on the southwestern coast of Pakistan. Gwadar started operation in March 2008, with 3 multi-purpose berths, a 602-meter quay, and 12.5-meter depth. Another 9 berths are under construction, with a 20-meter depth. In 2015, it was announced that the port would be leased to the Chinese till 2059 and further developed under the China-Pakistan Economic Corridor.
Terminals
| Technology | Specific piers and ports | null |
1230089 | https://en.wikipedia.org/wiki/Euoplocephalus | Euoplocephalus | Euoplocephalus ( ) is a genus of large herbivorous ankylosaurid dinosaurs, living during the Late Cretaceous of Canada. It has only one named species, Euoplocephalus tutus.
The first fossil of Euoplocephalus was found in 1897 in Alberta. In 1902, it was named Stereocephalus, but that name had already been given to an insect, so it was changed in 1910. Later, many more ankylosaurid remains were found from the Campanian of North America and often made separate genera. In 1971, Walter Coombs concluded that they all belonged to Euoplocephalus, which then make it one of the best-known dinosaurs. Recently, however, experts have come to the opposite conclusion, limiting the authentic finds of Euoplocephalus to about a dozen specimens. These include a number of almost complete skeletons, so nevertheless much is known about the build of the animal.
Euoplocephalus reached in length and in body mass. Its body was low-slung and very flat and wide, standing on four sturdy legs. Its head had a short drooping snout with a horny beak to bite off plants that were digested in the large gut. Like other ankylosaurids, Euoplocephalus was largely covered by bony armor plates, among them rows of large high-ridged oval scutes. The neck was protected by two bone rings. It could also actively defend itself against predators like Albertosaurus, Daspletosaurus, and Gorgosaurus using a heavy club at the end of its tail.
Description
Size
Among the ankylosaurids, Euoplocephalus was exceeded in size only by Ankylosaurus, and perhaps Tarchia and Cedarpelta. Euoplocephalus was about long and weighed about . Like other ankylosaurids, it had a very broad and flat low-slung torso, about four feet high, positioned on four short legs.
Distinguishing traits
The skull of Euoplocephalus can be distinguished from most other ankylosaurids by several anatomical details, including: the pattern of bony sculpturing in the region in front of the eyes; the form of the palpebral bones (small bones over the eyes), which may have served as bony eyelids; the shallowness of the nasal vestibule at the entrance of the nasal cavity; the medial curve of the tooth rows in the upper jaw; and the teeth, which are relatively small, lacking true cingula, and having variable fluting of the denticles. However, these traits are shared with a number of closely related forms, some of which in the past have been included in the genus. Combining such forms, Walter Coombs and Teresa Maryańska in 1990 stated that Euoplocephalus could be distinguished based on four traits. The premaxillae, the front snout bones, are not covered by dermal ossifications. The external bony nostrils are slit-like, face to the front and are each divided by a vertical bone sheet or septum. The beak has a width equal to or greater than the distance between the rear maxillary, upper cheek, teeth rows. The foot has three digits, metatarsals with toes.
In 2013, Victoria Arbour and Phil Currie provided a differential diagnosis, setting Euoplocephalus apart from its nearest relatives. When compared with Anodontosaurus and Scolosaurus, Euoplocephalus lacks round osteoderms at the base of the squamosal and quadratojugal horns. Compared with Anodontosaurus it lacks small osteoderms at the bases of the large osteoderms on the first cervical half-ring, but, contrary to that genus, does in top view have large rounded osteoderms at the sides of the tail club. It differs from Dyoplosaurus in possessing sacral ribs that perpendicularly point outwards. It differs from Scolosaurus in possessing keeled osteoderms with a round or oval base on the top and sides of the first cervical half-ring and having a shorter rear blade of the ilium. Euoplocephalus differs from Ankylosaurus in possessing anteriorly directed external nostrils and in lacking a continuous keel between the squamosal horn and the supraorbitals.
Skeleton
The skull of Euoplocephalus resembles a truncated equilateral triangle when viewed from above, and is slightly wider than it is long. The largest known skull, that of specimen AMNH 5403, has a length of 411 millimetres and a width of 478 millimetres. Most skull elements are completely fused and two skull openings normally present with dinosaurs, the antorbital fenestra and the upper temporal fenestra, have closed. The skull has nineteen to twenty-four teeth in each upper jaw. The frontmost snout bone, the premaxilla, is toothless. The teeth are very small, with a maximal height and width of just 7.5 millimetres. The strongly drooping snout is blunt, wide and high, and filled with very complex air passages and sinuses, the form and function of which are not yet completely understood. Each side has two external nostrils. The lower jaw has a very concave upper rim with twenty-one teeth. At its front a short low extension is present, to provide contact with the wide predentary, the bony core of the lower horny beak, that fitted within the upper beak of the snout.
As in most quadrupedal ornithischians, its neck is moderately short. The scapula is massive and robust, and the very robust forelimbs are shorter than the hindlimbs. The tail is long and ends in a bony club. Old restorations of Euoplocephalus and rejected synonyms (Dyoplosaurus, Scolosaurus) often show a club with two large vertical spikes. This is an error based on a restoration of Scolosaurus by Franz Nopcsa; the specimen he used had an incomplete tail that stopped just beyond the pair of conical spikes now known to have been positioned halfway along its length. He restored the tail as ending just after the spines. Other artists combined the spikes with the tail club, compounding the inaccuracy. The narrow distal half of the tail is stiffened by bundles of ossified tendons.
The vertebral column of Euoplocephalus is made up of at least seven cervical vertebrae, at least eleven "free" dorsal vertebrae, typically four sacrodorsals forming a fused "sacral rod" in front of the sacrum proper, three "true" sacral vertebrae, and between one and four caudosacrals. Like with other ankylosaurians, the last four dorsal vertebrae and the first caudal vertebra are thus fused to the sacrum, forming a reinforced synsacrum of at least eight vertebrae, the spines of which form a fused supraneural plate, also incorporating the zygapophyses. There are at least twenty-one caudal vertebrae; the total number of caudal vertebrae is uncertain because approximately ten are fused to form part of the tail club, bringing the total to about thirty. This fusion is also seen in other ankylosaurids; it is possible that the extent of fusion is an age-related feature.
The humerus is very robust with strongly expanded upper and lower joints, combined with a narrow shaft. On the upper shaft an enormous deltopectoral crest is present of which the lower part does not gradually merge with the shaft but is warped to the front, forming a thick knob or lip. All this indicates a very heavy musculature. In the lower arm the robust ulna has a well-developed olecranon process. The wrist and hand bones are not well known. In the pelvis, the front blade of the ilium splayed out to the front, reaching all the way to the widest point of the belly to support the gut. This blade also forms a bone shelf at the rear side of the body. The rear blade of the ilium is shorter than the diameter of the hip socket it was located behind, meaning the leg is located at the rear end of the pelvis, near the tail base and much closer to the midline than the belly sides. The pubic bone is unknown. The ischium is a short, curved, vertically positioned bone strap. The thighbone is short, robust and straight with a low fourth trochanter positioned below the midpoint of the shaft. The robust shinbone is shorter than the thighbone. The foot is not well known but functionally tridactyl with hoof-shaped instead of sharp claws.
Armor
The head and body of Euoplocephalus were covered with bony armor, except for parts of the limbs and possibly the distal tail. The armor consisted of osteoderms, skin ossifications that are not part of the skeleton proper. This armor was in 1982 extensively described by Kenneth Carpenter, who however, largely based himself on the very complete specimen NHMUK R5161, the holotype of Scolosaurus, which genus no longer is seen as a synonym of Euoplocephalus. When limited to the certain material of the latter genus, little is known of the exact configuration of the armor, with the exception of the head and the neck. The most informative specimen in this respect would then be ROM 1930, having conserved some osteoderms of the torso in their original position.
In any case, much of the armor was made up of small ossicles, bony round scutes with a diameter of less than five millimetres, of which often hundreds have been found with a single specimen. If the armor was configured in an identical way to that of Scolosaurus, many of these small ossicles had fused into a kind of pavement, forming transverse bands on the body. The banded arrangement is thought to have permitted some freedom of movement. Four of these bands might have been present on the anterior half of the tail, three on the pelvis, perhaps fused into a single "sacral shield", and four across the front part of the torso. Inset in these bands were horizontal rows of larger oval, flat or keeled, scutes. Types of large scutes varied by body region. It might be that the scutes on the shoulder, near the midline of the body, were largest and tallest; ROM 1930 includes some osteoderms with a base length of fifteen centimeters. Little is known about the armor of the limbs. Large keeled plates were present on the upper arms as shown by specimen TMP 1997.132.01 conserving a round osteoderm near the humerus with a diameter of twenty centimeters and narrower spikes associated with the lower arm. The neck was protected by two bone rings, open at the underside, that are called "cervical half-rings". Earlier seen as a fusion of osteoderms, this was doubted by Arbour et al. in 2013, who pointed out that they formed a lower layer, possibly consisting of ossified cartilage, as indicated by a smooth surface and a woven bone texture. Each half-ring is constructed out of six rectangular concave plates, three per side. Each plate has a large keeled osteoderm on top, often not fused with it. With Euoplocephalus, these neck osteoderms do not have smaller osteoderms at their bases, and their keels do not overhang their posterior edges.
The armor of the skull consists of a large number of ossicles, called caputegulae ("head tiles"), that have fused with the normal skull elements, largely fading their sutures. On the snout they form a chaotic and asymmetric mosaic. On the rear nasal region, at the midline a single hexagonal larger plate is present. A keeled plate per side, somewhat more to the back forms the snout rim. Behind the level of the eye sockets the caputegulae fuse into a single bone surface. The upper rim of the eye sockets is formed by two pyramid-shaped osteoderms pointing to the sides and rear. In addition, Euoplocephalus had two pyramid-shaped squamosal "horns" growing from the back corners of its head. Between them the nuchal crest is covered by two osteoderms per side. At the lower rear side of the skull, a quadratojugal horn is present, in the form of an enormous tongue-shaped osteoderm projecting below.
Discovery and species
Canadian paleontologist Lawrence Morris Lambe discovered the first specimen on 18 August 1897 in the area of the present Dinosaur Provincial Park, in the valley of the Red Deer River, Alberta, Canada. In 1902, this fossil, CMN 210 (also NMC 210) was designated as the holotype specimen of the type species Stereocephalus tutus. This specimen consists of the upper part of a cranium and a transverse series of five scutes that were part of a cervical half ring. The generic name was derived from Greek στερεός, stereos, "solid", and κεφαλή, kephalè, "head", which refers to the formidable armour. However, the genus name was already preoccupied — the name had already been given to an insect, the beetle Stereocephalus Lynch 1884 — so Lambe changed it to Euoplocephalus in 1910, with as combinatio nova (new combination name) Euoplocephalus tutus. The type species remains Stereocephalus tutus. In 1915, Edwin Hennig classified E. tutus under the genus Palaeoscincus Leidy 1856, coining a Palaeoscincus tutus. Today however, Palaeoscincus is considered to be a nomen dubium based on indeterminate ankylosaurian teeth. In 1964, Euoplocephalus was by Oskar Kuhn referred to Ankylosaurus, as a Ankylosaurus tutus.
The genus name Euoplocephalus, meaning "well-armed head", is derived from the Greek words eu (εὖ) meaning "well", hoplo~ (ὁπλο~) meaning "armed", and kephale (κεφαλή) meaning "head". This name has been misspelled more than a dozen different ways in formal scientific literature. The specific name tutus means "safely protected" in Latin. The only valid species known today is Euoplocephalus tutus.
Referred material
During the early twentieth century many more ankylosaurid fossils were uncovered in North America. Some were referred to Euoplocephalus, others named as separate genera. In 1971 however, Walter Coombs submitted a dissertation containing a landmark re-appraisal of North American ankylosaurs. He noted that, among the many specimens similar to Euoplocephalus, their skulls varied so much that either every known specimen must be a new species, or they all represented individual variation within a single species: Euoplocephalus tutus. Starting from this assumption that there was only one species of ankylosaur during the Campanian stage of the Upper Cretaceous, Coombs synonymized the genera Anodontosaurus, Dyoplosaurus, and Scolosaurus with Euoplocephalus and the species A. lambei, D. acutosquameus, and S. cutleri with E. tutus, creating a species that spanned nearly ten million years, or the entire Campanian. The fossils now referred to this species contained more than forty individuals discovered in Alberta, Canada and Montana in the United States, which would have made Euoplocephalus the best known ankylosaurid. This included fifteen skulls, teeth, and a few almost-complete skeletons, found with the armor still attached. Individual armor plates are the most commonly found element from them. In 1978, Coombs even included the Asian ankylosaurid Tarchia in the genus, renaming it as Euoplocephalus giganteus.
The synonymy of all Campanian North-American ankylosaurids was followed for several decades, until scientists from the University of Alberta began to re-examine the fossils. A 2009 study found that Dyoplosaurus is in fact a valid taxon, and identified unique characteristics that differentiated it from Euoplocephalus, including its triangular claws. Victoria Arbour (2010) argued that Anodontosaurus (known from the Horseshoe Canyon Formation) is distinct from Euoplocephalus and is also a valid taxon; according to Arbour, Anodontosaurus differs from Euoplocephalus in distinctive skull and cervical half ring ornamentation, as well as tail club morphology, including the presence of pointed, triangular knob osteoderms in Anodontosaurus. Furthermore, Arbour (2010) suggested to reassign all Horseshoe Canyon Formation ankylosaurine specimens from Euoplocephalus to Anodontosaurus. The validity of Anodontosaurus was accepted in two subsequent studies. The first, published by Paul Penkalski and William T. Blows in 2013, re-validated Scolosaurus as well. The second study, by Penkalski (2013), named and described Oohkotokia from Montana on the basis of remains that were originally thought to be referable to Euoplocephalus.
Palaeoscincus asper, "the rough one", is now considered to be Euoplocephalus. It is a dubious tooth taxon from the late Campanian Dinosaur Park Formation of Alberta, named by Lambe in 1902. It consists of a single tooth, specimen NMC 1349.
In 2013 Arbour limited the specimens that could be reliably referred to Euoplocephalus to the lowest thirty metres of the Dinosaur Park Formation. The material would in that case, apart from the holotype, consist of partial skeletons with skull AMNH 5337, AMNH 5403, AMNH 5404, AMNH 5405, ROM 1930 and UALVP 31; partial skeleton lacking the skull AMNH 5406; CMN 842, a cervical half-ring; CMN 8876, a skull, TMP 1979.14.74, a fragmentary skull; and UALVP 47977, a skull roof piece. The hands, feet and tail, including the club, are therefore not completely known. However, many of those specimens have now been reassigned to other, new taxa, including Scolosaurus and Platypelta, and others.
Classification
In 1910, Lambe assigned Euoplocephalus to the Stegosauria, a group then encompassing all armoured dinosaur forms and thus having a much wider range than the present concept. In 1917, Charles Whitney Gilmore assigned it to the Ankylosauridae. Today, Euoplocephalus is still seen as an ankylosaurid, but as a member of the Ankylosauria, not the Stegosauria. It is likely also a member of the derived subgroup Ankylosaurinae. The recent splitting of the ankylosaurid Campanian material of North America has complicated the issue of the direct affinities of Euoplocephalus. Penkalski (2013) performed a small phylogenetic analysis of some ankylosaurine specimens. The only Anodontosaurus specimen that was included in this analysis was its holotype. Anodontosaurus was placed in a polytomy with the holotype of Euoplocephalus and some specimens that are referred to it, while Oohkotokia was placed in a clade with Dyoplosaurus, and specimens that are thought to represent either Dyoplosaurus or Scolosaurus.
The following cladogram is based on a 2015 phylogenetic analysis of the Ankylosaurinae conducted by Arbour and Currie:
The results of an earlier analysis of the ankylosaurid tree by Thompson et al. (2011), is shown by this cladogram.
Paleoecology
Euoplocephalus, following the synonymizations proposed by Coombs (1971), was thought to exist for far longer, and was a member of more distinct faunas, than any of its contemporaries, as these fossils date to between 76.5 and 67 million years ago, in the Campanian-Maastrichtian ages of the late Cretaceous period, and came from the Dinosaur Park and Horseshoe Canyon Formations of Alberta, Two Medicine Formation of Montana, and possibly from the Oldman Formation of Alberta. Fossils that were initially believed to be from the Judith River Formation of Montana, are actually from the Dinosaur Park Formation. However, recent studies referred all Horseshoe Canyon Formation specimens to Anodontosaurus, and all Two Medicine Formation specimens to Oohkotokia. A specimen from the lowermost Dinosaur Park, or possibly from the underlying Oldman Formation, was reassigned back to Scolosaurus. Although the stratigraphic range of the holotype of Euoplocephalus is uncertain, all specimens that can be reliably referred to E. tutus came from the lower 40 m and the upper >10 m of the Dinosaur Park Formation. There are no known ankylosaurids from the top 20–25 m of the Formation. Thus, all Euoplocephalus specimens date to between 76.4 (or less) and 75.6 million years ago, in the late Campanian stage.
Paleobiology
According to Coombs, Euoplocephalus may have had cursorial abilities on par with those of the modern rhinoceros and hippopotamus. Based on the form of the humerus-shoulder articulation and the arrangement of the protracting muscles of the upper arm, it appears that the upper arm sloped away from the body. Coombs and Maryanska (1990) observed that Euoplocephalus specimens are usually discovered as isolated elements or partial skeletons, which suggested that this animal engaged in solitary habits and was usually either solitary or participated in small group clusters.
The armor of Euoplocephalus may have had a keratinous covering, or it may have floated in the skin, as is seen in modern crocodiles. In addition to protection, the heavily vascularized armor may have had a role in thermoregulation. The palpebral bones over the eyes may have provided additional protection for the eyes. Such bones with Euoplocephalus have been discovered in the upper part of the eye socket, instead of in front of the upper socket rim which is the more common position. Coombs explained this by assuming that these bones were located in the eyelid musculature and were probably mobile enough to be moved over the eyes.
Defense
The tail club of ankylosaurids has often been interpreted as a defensive weapon. In Euoplocephalus, the presence of ossified tendons only with the distal half of the tail may support such a function. Because only the distal half of the tail was stiffened by tendons, the anterior half could still move freely from side to side. The ossified tendons would have transmitted the force of the swing to the club and reinforced the supporting vertebrae. The club was likely held just above the ground, as there was not sufficient musculature to raise the tail very high. A 2009 study concluded that "large ankylosaurian clubs could generate sufficient force to break bone during impacts, while average and small ones could not". It has also concluded that "tail swinging behavior is feasible in ankylosaurids, but it remains unknown whether the tail was used for interspecific defense, intraspecific combat, or both". The tail club could be swung low, toward the fragile metatarsals or shin bones of attacking theropods.
Senses and airflow
Euoplocephalus had relatively small eyes, but this does not necessarily mean that it had restricted vision. The complex respiratory passages observed in the skull suggest that Euoplocephalus had a good sense of smell, although in 1978 an examination of casts of the endocranium did not show an enlarged olfactory region of the brain. Teresa Maryanska, who has worked extensively on Mongolian ankylosaurids, suggested that the respiratory passages were primarily used to perform a mammal-like treatment of inhaled air, based on the presence and arrangement of specialized bones, which are present in Euoplocephalus.
A 2011 study found that the nasal passages of Euoplocephalus were looped and complex; possibly an adaptation for heat and water balance and vocal resonance, and researchers discovered an enlarged and vascularised chamber at the back of the nasal tract, which was considered by the authors to be an adaptation to improve its sense of smell. The researchers also managed to reconstruct the dinosaur's inner ear and concluded that it was capable of hearing at low frequencies. They suggested that this may have been an adaption to hearing low-toned resonant sounds produced by the nasal passages.
Diet
Euoplocephalus, like other ankylosaurians, is thought to have been a herbivore. It had a broad muzzle, which could indicate that it was a non-selective feeder, perhaps similar to a hippopotamus. This would provide niche separation from contemporaneous nodosaurids with narrower muzzles. Ankylosaurians have historically been thought of as feeding using simple up-and-down movements of the jaws. Georg Haas (1969) examined the evidence for the jaw muscles of two skulls (AMNH 5337 and 5405) and concluded that despite the large size of the skulls the associated musculature was relatively weak. He also thought jaw movement was largely orthogonal, in the vertical plane only. Haas extrapolated from this that dinosaurs like Euoplocephalus likely ate relatively soft non-abrasive vegetation.
However, later research indicated that forward and sideways jaw movement was possible, the skull being able to withstand considerable forces. Euoplocephalus appears to have been able to make more complex movements. Tooth wear and jaw articulations (within the lower jaw and at the lower jaw-quadrate joint) suggest that the lower jaws were pulled back during feeding, and also slightly pivoted inward. This action would have sheared food. A study conducted in 2014 found that the ankylosaurs were capable of eating of tough fibrous plant material, though not to the same degree as their nodosaur relatives or the ceratopsians and hadrosaurs.
| Biology and health sciences | Ornitischians | Animals |
23582155 | https://en.wikipedia.org/wiki/Grid%20code | Grid code | A grid code is a technical specification which defines the parameters a facility connected to a public electric grid has to meet to ensure safe, secure and economic proper functioning of the electric system. The facility can be an electricity generating plant, a consumer, or another network.
The grid code is specified by an authority responsible for the system integrity and network operation.
Its elaboration usually implicates network operators (distribution or transmission system operators), representatives of users and, to an extent varying between countries, the regulating body.
Contents of a grid code vary depending on the transmission company's requirements. Typically, a grid code will specify the required behavior of a connected generator during system disturbances. These include voltage regulation, power factor limits and reactive power supply, response to a system fault (e.g. short-circuit), response to frequency changes on the grid, and requirement to "ride through" short interruptions of the connection.
There is not a common grid code in all countries and each electric grid has its own grid code. Even in North America, there is no grid code that applies to all territories.
Independent power producers
All generators including Independent power producers like photovoltaic power stations or wind farms have to comply with the grid code.
Categorizing
Grid code requirements can be divided into two categories: static and dynamic requirements.
| Technology | Concepts | null |
764405 | https://en.wikipedia.org/wiki/Turing%20degree | Turing degree | In computer science and mathematical logic the Turing degree (named after Alan Turing) or degree of unsolvability of a set of natural numbers measures the level of algorithmic unsolvability of the set.
Overview
The concept of Turing degree is fundamental in computability theory, where sets of natural numbers are often regarded as decision problems. The Turing degree of a set is a measure of how difficult it is to solve the decision problem associated with the set, that is, to determine whether an arbitrary number is in the given set.
Two sets are Turing equivalent if they have the same level of unsolvability; each Turing degree is a collection of Turing equivalent sets, so that two sets are in different Turing degrees exactly when they are not Turing equivalent. Furthermore, the Turing degrees are partially ordered, so that if the Turing degree of a set X is less than the Turing degree of a set Y, then any (possibly noncomputable) procedure that correctly decides whether numbers are in Y can be effectively converted to a procedure that correctly decides whether numbers are in X. It is in this sense that the Turing degree of a set corresponds to its level of algorithmic unsolvability.
The Turing degrees were introduced by and many fundamental results were established by . The Turing degrees have been an area of intense research since then. Many proofs in the area make use of a proof technique known as the priority method.
Turing equivalence
For the rest of this article, the word set will refer to a set of natural numbers. A set X is said to be Turing reducible to a set Y if there is an oracle Turing machine that decides membership in X when given an oracle for membership in Y. The notation X ≤T Y indicates that X is Turing reducible to Y.
Two sets X and Y are defined to be Turing equivalent if X is Turing reducible to Y and Y is Turing reducible to X. The notation X ≡T Y indicates that X and Y are Turing equivalent. The relation ≡T can be seen to be an equivalence relation, which means that for all sets X, Y, and Z:
X ≡T X
X ≡T Y implies Y ≡T X
If X ≡T Y and Y ≡T Z then X ≡T Z.
A Turing degree is an equivalence class of the relation ≡T. The notation [X] denotes the equivalence class containing a set X. The entire collection of Turing degrees is denoted .
The Turing degrees have a partial order ≤ defined so that [X] ≤ [Y] if and only if X ≤T Y. There is a unique Turing degree containing all the computable sets, and this degree is less than every other degree. It is denoted 0 (zero) because it is the least element of the poset . (It is common to use boldface notation for Turing degrees, in order to distinguish them from sets. When no confusion can occur, such as with [X], the boldface is not necessary.)
For any sets X and Y, X join Y, written X ⊕ Y, is defined to be the union of the sets } and }. The Turing degree of X ⊕ Y is the least upper bound of the degrees of X and Y. Thus is a join-semilattice. The least upper bound of degrees a and b is denoted a ∪ b. It is known that is not a lattice, as there are pairs of degrees with no greatest lower bound.
For any set X the notation X′ denotes the set of indices of oracle machines that halt (when given their index as input) when using X as an oracle. The set X′ is called the Turing jump of X. The Turing jump of a degree [X] is defined to be the degree [X′]; this is a valid definition because X′ ≡T Y′ whenever X ≡T Y. A key example is 0′, the degree of the halting problem.
Basic properties of the Turing degrees
Every Turing degree is countably infinite, that is, it contains exactly sets.
There are distinct Turing degrees.
For each degree a the strict inequality a < a′ holds.
For each degree a, the set of degrees below a is countable. The set of degrees greater than a has size .
Structure of the Turing degrees
A great deal of research has been conducted into the structure of the Turing degrees. The following survey lists only some of the many known results. One general conclusion that can be drawn from the research is that the structure of the Turing degrees is extremely complicated.
Order properties
There are minimal degrees. A degree a is minimal if a is nonzero and there is no degree between 0 and a. Thus the order relation on the degrees is not a dense order.
The Turing degrees are not linearly ordered by ≤T.
In fact, for every nonzero degree a there is a degree b incomparable with a.
There is a set of pairwise incomparable Turing degrees.
There are pairs of degrees with no greatest lower bound. Thus is not a lattice.
Every countable partially ordered set can be embedded in the Turing degrees.
An infinite strictly increasing sequence a1, a2, ... of Turing degrees cannot have a least upper bound, but it always has an exact pair c, d such that ∀e (e<c∧e<d ⇔ ∃i e≤ai), and thus it has (non-unique) minimal upper bounds.
Assuming the axiom of constructibility, it can be shown there is a maximal chain of degrees of order type .
Properties involving the jump
For every degree a there is a degree strictly between a and a′. In fact, there is a countable family of pairwise incomparable degrees between a and a′.
Jump inversion: a degree a is of the form b′ if and only if 0′ ≤ a.
For any degree a there is a degree b such that a < b and b′ = a′; such a degree b is called low relative to a.
There is an infinite sequence ai of degrees such that a′i+1 ≤ ai for each i.
Post's theorem establishes a close correspondence between the arithmetical hierarchy and finitely iterated Turing jumps of the empty set.
Logical properties
showed that the first-order theory of in the language or is many-one equivalent to the theory of true second-order arithmetic. This indicates that the structure of is extremely complicated.
showed that the jump operator is definable in the first-order structure of with the language .
Recursively enumerable Turing degrees
A degree is called recursively enumerable (r.e.) or computably enumerable (c.e.) if it contains a recursively enumerable set. Every r.e. degree is below 0′, but not every degree below 0′ is r.e.. However, a set is many-one reducible to 0′ iff is r.e..
: The r.e. degrees are dense; between any two r.e. degrees there is a third r.e. degree.
and : There are two r.e. degrees with no greatest lower bound in the r.e. degrees.
and : There is a pair of nonzero r.e. degrees whose greatest lower bound is 0.
: There is no pair of r.e. degrees whose greatest lower bound is 0 and whose least upper bound is 0′. This result is informally called the nondiamond theorem.
: Every finite distributive lattice can be embedded into the r.e. degrees. In fact, the countable atomless Boolean algebra can be embedded in a manner that preserves suprema and infima.
: Not all finite lattices can be embedded in the r.e. degrees (via an embedding that preserves suprema and infima). A particular example is shown to the right. L. A. Harrington and T. A. Slaman (see ): The first-order theory of the r.e. degrees in the language 〈 0, ≤, = 〉 is many-one equivalent to the theory of true first-order arithmetic.
Additionally, there is Shoenfield's limit lemma, a set A satisfies iff there is a "recursive approximation" to its characteristic function: a function g such that for sufficiently large s, .
A set A is called n-r e. if there is a family of functions such that:
As is a recursive approximation of A: for some t, for any s≥t we have As(x) = A(x), in particular conflating A with its characteristic function. (Removing this condition yields a definition of A being "weakly n-r.e.")
As is an "n-trial predicate": for all x, A0(x)=0 and the cardinality of is ≤n.
Properties of n-r.e. degrees:
The class of sets of n-r.e. degree is a strict subclass of the class of sets of (n+1)-r.e. degree.
For all n>1 there are two (n+1)-r.e. degrees a, b with , such that the segment contains no n-r.e. degrees.
and are (n+1)-r.e. iff both sets are weakly-n-r.e.
Post's problem and the priority method
Emil Post studied the r.e. Turing degrees and asked whether there is any r.e. degree strictly between 0 and 0′. The problem of constructing such a degree (or showing that none exist) became known as Post's problem. This problem was solved independently by Friedberg and Muchnik in the 1950s, who showed that these intermediate r.e. degrees do exist (Friedberg–Muchnik theorem). Their proofs each developed the same new method for constructing r.e. degrees, which came to be known as the priority method. The priority method is now the main technique for establishing results about r.e. sets.
The idea of the priority method for constructing a r.e. set X is to list a countable sequence of requirements that X must satisfy. For example, to construct a r.e. set X between 0 and 0′ it is enough to satisfy the requirements Ae and Be for each natural number e, where Ae requires that the oracle machine with index e does not compute 0′ from X and Be requires that the Turing machine with index e (and no oracle) does not compute X. These requirements are put into a priority ordering, which is an explicit bijection of the requirements and the natural numbers. The proof proceeds inductively with one stage for each natural number; these stages can be thought of as steps of time during which the set X is enumerated. At each stage, numbers may be put into X or forever (if not injured) prevented from entering X in an attempt to satisfy requirements (that is, force them to hold once all of X has been enumerated). Sometimes, a number can be enumerated into X to satisfy one requirement but doing this would cause a previously satisfied requirement to become unsatisfied (that is, to be injured). The priority order on requirements is used to determine which requirement to satisfy in this case. The informal idea is that if a requirement is injured then it will eventually stop being injured after all higher priority requirements have stopped being injured, although not every priority argument has this property. An argument must be made that the overall set X is r.e. and satisfies all the requirements. Priority arguments can be used to prove many facts about r.e. sets; the requirements used and the manner in which they are satisfied must be carefully chosen to produce the required result.
For example, a simple (and hence noncomputable r.e.) low X (low means X′=0′) can be constructed in infinitely many stages as follows. At the start of stage n, let Tn be the output (binary) tape, identified with the set of cell indices where we placed 1 so far (so X=∪n Tn; T0=∅); and let Pn(m) be the priority for not outputting 1 at location m; P0(m)=∞. At stage n, if possible (otherwise do nothing in the stage), pick the least i<n such that ∀m Pn(m)≠i and Turing machine i halts in <n steps on some input S⊇Tn with ∀m∈S\Tn Pn(m)≥i. Choose any such (finite) S, set Tn+1=S, and for every cell m visited by machine i on S, set Pn+1(m) = min(i, Pn(m)), and set all priorities >i to ∞, and then set one priority ∞ cell (any will do) not in S to priority i. Essentially, we make machine i halt if we can do so without upsetting priorities <i, and then set priorities to prevent machines >i from disrupting the halt; all priorities are eventually constant.
To see that X is low, machine i halts on X iff it halts in <n steps on some Tn such that machines <i that halt on X do so <n-i steps (by recursion, this is uniformly computable from 0′). X is noncomputable since otherwise a Turing machine could halt on Y iff Y\X is nonempty, contradicting the construction since X excludes some priority i cells for arbitrarily large i; and X is simple because for each i the number of priority i cells is finite.
| Mathematics | Computability theory | null |
764520 | https://en.wikipedia.org/wiki/Fishing%20lure | Fishing lure | A fishing lure is any one of a broad category of artificial angling baits that are inedible replicas designed to mimic prey animals (e.g. baitfish, crustaceans, insects, worms, etc.) that attract the attention of predatory fish, typically via appearances, flashy colors, bright reflections, movements, vibrations and/or loud noises which appeal to the fish's predation instinct and entice it into gulping the lure. Angling activities using lures are known as lure fishing.
As a terminal tackle, lures are attached to the end of a fishing line and most are equipped with one or more hooks that come in various styles (most commonly treble hooks). They are designed to be fake foods that fool carnivorous fish into an aggressive strike, the force of which will alert the fisherman to yank the line and secure a hookset inside the fish mouth. However some hookless lures are also used merely to bait the fish nearer, so it can be hooked by another lure or be captured via other means such as netting, trapping, impaling with a spear or pole hook, snagging, shooting or even catching by hand.
Most lures are commercially made, but some are hand-made by the angler (such as hand-tied fly lures, which is considered a hobbyist challenge by many amateur entomologists). Modern lures are typically cast and manipulated with a fishing rod and retrieved with a reel, but there are some who use a technique where the line is directly held with hands (known as handlining) rather than manipulated through a rod. Longlining and trolling, commonly used for commercial fishing, also can employ lures to catch fish. When used for recreational fishing, the lure is repeatedly cast out far and then reeled back towards the angler, and during retrieval it interacts with the water current and drag, creating vibrations, turbulent splashes and/or a popping action. A skilled angler can explore many possible hiding places for fish through lure casting under logs, grass and on flats.
History
In early time, angling used predominantly edible baits impaled on hooks made from animal bones or bronze. The ancient Chinese and Egyptians practised angling with fishing rods, hooks and lines as early as 2,000 B.C., though most of the first fishermen used handlines. The Roman scholar Claudius Aelianus first described the practice of "fasten red wool... round a hook, and fit on to the wool two feathers which grow under a cock's wattles" by Macedonian anglers on the Astraeus River, an early form of fly fishing, near the end of the 2nd century. The Chinese were the first to make modern-looking fishing line, spun from fine silk, and the use of "wooden fish", likely an early type of surface lure, to catch larger fish had been wide spread at least since the Song dynasty (960–1279).
Nordic people have been making spoon lures from the 8th-13th century AD. Most of the lures are made from iron, bronze, copper, and in one case an iron hook soldered to a copper spoon. Many lures had varying shapes and sizes fitting different scenarios like ice fishing and summer fishing. Modern spoon lures appear to have originated in Scandinavia in the late 1700s.
English tackle shops are recorded as selling tin minnows in the middle of the 18th century, and realistic imitations of bugs and grubs made from painted rubber appeared as early as 1800. Early English minnow baits were largely designed to spin as their attracting action, as exemplified by the “Devon”-style lure first produced in quantity by F. Angel of Exeter. The number and variety of artificial baits increased dramatically in the mid- to late 19th century.
The first production lures made in the United States, mostly metal spoons and spinnerbaits, came on the market in the last half of the 19th century. The makers included Julio T. Buel, Riley Haskell, W. D. Chapman and Enterprise Manufacturing Company. Modern fishing plugs were first made commercially in the United States in the early 1900s by firms including Heddon in Michigan and Enterprise Mfg. (Pflueger) in Ohio. Before this time most fishing lures were made by individual craftsman. Commercial-made lures were based on the same ideas that the individual craftsmen were making but on a larger scale.
Methods
The fishing lure is either directly tied to a fishing line (usually a leader) by a knot such as the improved clinch knot or the Palomar knot, or linked to the line via a small split ring (which allows more freedom of motion) and/or a tiny safety pin-like fastener called a "snap", which is usually also connected to a swivel. The fishing line is in turn connected to a fishing reel which cranks in and releases out the line, and is manipulated by the fishing rod via a series of ring guides that impart lateral displacement on the line.
The in-water motion of the lure is generated by winding the line back towards the angler, by sweeping the fishing rod sideways, jigging movements with the rod tip, or by being towed behind a moving boat (trolling). These movements mimic the behaviors of small preys, which draw the attention of larger aquatic predators and fool them into identifying the lure as an easy meal. Exceptions include artificial flies, commonly just called flies by fly fishers, which either float on the water surface, slowly sink or float underwater, and represent some form of drowning insect.
Types
There are many types of fishing lures. Today's modern definition for lures are that they be made of wood, plastic, rubber, metal, cork, and materials like feathers, animal hair, string, tinsel and others. They could also have any number of moving parts or no moving parts. They can be retrieved fast or slow. Some of the lures can be used alone, or with another lure.
In most cases they are manufactured to resemble prey for the fish, but they are sometimes engineered to appeal to a fishes' sense of territory, curiosity or aggression. Most lures are made to look like dying, injured, or fast moving fish. They include the following types:
Artificial flies are designed to resemble all manner of insect prey and are used in fly fishing.
Combined lures combine properties of several different types of lures.
Chatterbait, also known as "bladed jig" or "vibrating jig", is an amalgamation of several lure constructs. It has a weighted hook (jighead), a feathered/silicone stranded skirt, and an inline spinner blade. As it is cast, and retrieved, it vibrates under water alerting nearby fish of a potential snack.
Fish decoy is a type of lure that traditionally was carved to resemble a fish, frog, small rodent, or an insect that lures in fish so they can be speared. They are often used through the ice by fishermen and also by Inuit as part of their diet. The Mitchell Museum of the American Indian collection includes Native American fish decoys. William Jesse Ramey is considered a vintage master carver of fish decoys, and his work has been featured in museums.
Jigs are a weighted hook with a lead head opposite the sharp tip. They usually have a minnow or crawfish or even a plastic worm on it to get the fish's attention. Deep water jigs used in saltwater fishing consist of a large metallic weight, which gives the impression of the body of the bait fish, which has a hook attached via a short length of kevlar usually to the top of the jig. Some jigs can be fished in water depths down to 300 meters.
LED lures have a built in LED and battery to attract fish. They use a flashing or sometimes strobing pattern, using a combination of colors and LEDs.
Plugs are also known as crankbaits or minnows. These lures look like fish and they are run through the water where they can move in different ways because of instability due to the bib at the front under the head.
Soft plastic baits are lures made of plastic or rubber designed to look like fish, crabs, squid, worms, lizards, frogs, leeches and other creatures.
Spinnerbait are pieces of wire that are bent at about a 60- to 90-degree angle with a hook at the bottom and a flashy spinner at the top.
Spoon lures usually look like a spoon, with a wide rounded end, catching water to force action, and a narrower pointed end at the knot, similar in shape to a concave spearhead. It is shaped to have its center line off center to force the water to act upon it. They flash in the light while wobbling and darting due to their shape, which attracts fish.
Surface lures are also known as top water lures, poppers and stickbaits. They float and look like fish prey that is on top of the water. They can make a popping, burbling, or even a buzzing sound. It takes a long time to learn how to use this lure effectively. There are specific techniques for using surface lures effectively like "walking" them which gives a natural swimming look.
Swimbait is a soft plastic or wooden bait/lure that resembles an actual bait fish. Some of these have a tail that makes the lure/bait look like it is swimming when drawn through the water. Such a one made of wood would be hinged in certain places depending on its size.
One advantage of use of lure in general is the reduction in the use of live bait. This contributes to resolving one of the marine environment's more pressing problems; the undermining of marine food webs by overharvesting bait species which tend to occur lower in the food chain. Another advantage of lures is that their use promotes improved survival of fish during catch and release fishing. This is because lures reduce the incidence of deep hooking which has been correlated to fish mortality in many studies.
Rigs
A rig is an established terminal tackle setup that combines at least one hooked lure with one or more line sections, sinkers, bobbers, swivels, decorative beads, and sometimes other lures. A rig might be held by a rod, by hand, or attached to a boat or pier. Some rigs are designed to float near the surface of the water, others are designed to sink to the bottom. Some rigs are designed for trolling. Many rigs are designed especially for catching a single species of fish, but will work well for many different species.
Daisy chain
A daisy chain is a teaser rig consisting of a "chain" or cluster of plastic lures run without hooks, which mimics a school of forage fish that presents abundant food for predators. The purpose of a daisy chain is to attract pelagic fish to the stern of a boat into the lure "spread", which consists of a number of lures rigged with hooks.
Typically, a daisy chain's mainline is clear monofilament line with crimped-on droppers that connect the lure to the mainline. The last lure can be rigged with a hook or unrigged. The unrigged versions are used as teasers while the hooked versions are connected to a rod and reel. The lures used on a daisy chain are made from cedar plugs, plastic squids, jets, and other soft and/or hard plastic lures.
In some countries (e.g. New Zealand), daisy chains can sometimes refer to a rig which is used to catch baitfish in a similar arrangement to a "flasher rig" or a "sabiki rig"; a series of hooks with a small piece of colourful material/feather/plastic attached to each hook.
| Technology | Hunting and fishing | null |
765457 | https://en.wikipedia.org/wiki/Diethyl%20ether | Diethyl ether | Diethyl ether, or simply ether, is an organic compound with the chemical formula , sometimes abbreviated as . It is a colourless, highly volatile, sweet-smelling ("ethereal odour"), extremely flammable liquid. It belongs to the ether class of organic compounds. It is a common solvent. It was formerly used as a general anesthetic.
Production
Most diethyl ether is produced as a byproduct of the vapor-phase hydration of ethylene to make ethanol. This process uses solid-supported phosphoric acid catalysts and can be adjusted to make more ether if the need arises: Vapor-phase dehydration of ethanol over some alumina catalysts can give diethyl ether yields of up to 95%.
Diethyl ether can be prepared both in laboratories and on an industrial scale by the acid ether synthesis.
Uses
The dominant use of diethyl ether is as a solvent. One particular application is in the production of cellulose plastics such as cellulose acetate.
Laboratory solvent
It is a common solvent for the Grignard reaction in addition to other reactions involving organometallic reagents. These uses exploit its basicity. Diethyl ether is a popular non-polar solvent in liquid-liquid extraction. As an extractant, it is immiscible with and less dense than water.
Although immiscible, it has significant solubility in water (6.05 g/(100 ml) at 25 °C) and dissolves 1.5 g/(100 g) (1.0 g/(100 ml)) water at 25 °C.
Fuel
Diethyl ether has a high cetane number of 85–96 and, in combination with petroleum distillates for gasoline and diesel engines, is used as a starting fluid because of its high volatility and low flash point. Ether starting fluid is sold and used in countries with cold climates, as it can help with cold starting an engine at sub-zero temperatures. For the same reason it is also used as a component of the fuel mixture for carbureted compression ignition model engines.
Chemical reactions
Triethyloxonium tetrafluoroborate is prepared from boron trifluoride, diethyl ether, and epichlorohydrin:
Diethyl ether is a common laboratory aprotic solvent. It is susceptible to the formation of hydroperoxides.
Metabolism
A cytochrome P450 enzyme is proposed to metabolize diethyl ether.
Diethyl ether inhibits alcohol dehydrogenase, and thus slows the metabolism of ethanol. It also inhibits metabolism of other drugs requiring oxidative metabolism.
For example, diazepam requires hepatic oxidization whereas its oxidized metabolite oxazepam does not.
Safety, stability, regulations
Diethyl ether is extremely flammable and may form explosive vapour/air mixtures.
Since ether is heavier than air it can collect low to the ground and the vapour may travel considerable distances to ignition sources. Ether will ignite if exposed to an open flame, though due to its high flammability, an open flame is not required for ignition. Other possible ignition sources include – but are not limited to – hot plates, steam pipes, heaters, and electrical arcs created by switches or outlets. Vapour may also be ignited by the static electricity which can build up when ether is being poured from one vessel into another. The autoignition temperature of diethyl ether is . The diffusion of diethyl ether in air is (298 K, 101.325 kPa).
Ether is sensitive to light and air, tending to form explosive peroxides. Ether peroxides have a higher boiling point than ether and are contact explosives when dry. Commercial diethyl ether is typically supplied with trace amounts of the antioxidant butylated hydroxytoluene (BHT), which reduces the formation of peroxides. Storage over sodium hydroxide precipitates the intermediate ether hydroperoxides. Water and peroxides can be removed by either distillation from sodium and benzophenone, or by passing through a column of activated alumina.
Due to its application in the manufacturing of illicit substances, it is listed in the Table II precursor under the United Nations Convention Against Illicit Traffic in Narcotic Drugs and Psychotropic Substances as well as substances such as acetone, toluene and sulfuric acid.
History
The compound may have been synthesised by either Jābir ibn Hayyān in the 8th century or Ramon Llull in 1275. It was synthesised in 1540 by Valerius Cordus, who called it "sweet oil of vitriol" (oleum dulce vitrioli) – the name reflects the fact that it is obtained by distilling a mixture of ethanol and sulfuric acid (then known as oil of vitriol) – and noted some of its medicinal properties. At about the same time, Paracelsus discovered the analgesic properties of the molecule in dogs. The name ether was given to the substance in 1729 by August Sigmund Frobenius.
It was considered to be a sulfur compound until the idea was disproved in about 1800.
The synthesis of diethyl ether by a reaction between ethanol and sulfuric acid has been known since the 13th century.
Anesthesia
William T. G. Morton participated in a public demonstration of ether anesthesia on October 16, 1846, at the Ether Dome in Boston, Massachusetts. Morton had called his ether preparation, with aromatic oils to conceal its smell, "Letheon" after the Lethe River (Λήθη, meaning "forgetfulness, oblivion"). However, Crawford Williamson Long is now known to have demonstrated its use privately as a general anesthetic in surgery to officials in Georgia, as early as March 30, 1842, and Long publicly demonstrated ether's use as a surgical anesthetic on six occasions before the Boston demonstration. British doctors were aware of the anesthetic properties of ether as early as 1840 where it was widely prescribed in conjunction with opium. Diethyl ether was preferred by some practitioners over chloroform as a general anesthetic due to ether's more favorable therapeutic index, that is, a greater difference between an effective dose and a potentially toxic dose.
Diethyl ether does not depress the myocardium but rather it stimulates the sympathetic nervous system leading to hypertension and tachycardia. It is safely used in patients with shock as it preserves the baroreceptor reflex. Its minimal effect on myocardial depression and respiratory drive, as well as its low cost and high therapeutic index allows it to see continued use in developing countries. Diethyl ether could also be mixed with other anesthetic agents such as chloroform to make C.E. mixture, or chloroform and alcohol to make A.C.E. mixture. In the 21st century, ether is rarely used. The use of flammable ether was displaced by nonflammable fluorinated hydrocarbon anesthetics. Halothane was the first such anesthetic developed and other currently used inhaled anesthetics, such as isoflurane, desflurane, and sevoflurane, are halogenated ethers. Diethyl ether was found to have undesirable side effects, such as post-anesthetic nausea and vomiting. Modern anesthetic agents reduce these side effects.
Prior to 2005, it was on the World Health Organization's List of Essential Medicines for use as an anesthetic.
Medicine
Ether was once used in pharmaceutical formulations. A mixture of alcohol and ether, one part of diethyl ether and three parts of ethanol, was known as "Spirit of ether", Hoffman's Anodyne or Hoffman's Drops. In the United States this concoction was removed from the Pharmacopeia at some point prior to June 1917, as a study published by William Procter, Jr. in the American Journal of Pharmacy as early as 1852 showed that there were differences in formulation to be found between commercial manufacturers, between international pharmacopoeia, and from Hoffman's original recipe. It is also used to treat hiccups through instillation into the nasal cavity.
Recreational abuse
The recreational use of ether also took place at organised parties in the 19th century called ether frolics, where guests were encouraged to inhale therapeutic amounts of diethyl ether or nitrous oxide, producing a state of excitation. Long, as well as fellow dentists Horace Wells, William Edward Clarke and William T. G. Morton observed that during these gatherings, people would often experience minor injuries but appear to show no reaction to the injury, nor memory that it had happened, demonstrating ether's anaesthetic effects.
In the 19th century and early 20th century ether drinking was popular among Polish peasants. It is a traditional and still relatively popular recreational drug among Lemkos. It is usually consumed in a small quantity (kropka, or "dot") poured over milk, sugar water, or orange juice in a shot glass. As a drug, it has been known to cause psychological dependence, sometimes referred to as etheromania.
| Physical sciences | Esters and ethers | Chemistry |
765459 | https://en.wikipedia.org/wiki/Eurypterid | Eurypterid | Eurypterids, often informally called sea scorpions, are a group of extinct marine arthropods that form the order Eurypterida. The earliest known eurypterids date to the Darriwilian stage of the Ordovician period 467.3 million years ago. The group is likely to have appeared first either during the Early Ordovician or Late Cambrian period. With approximately 250 species, the Eurypterida is the most diverse Paleozoic chelicerate order. Following their appearance during the Ordovician, eurypterids became major components of marine faunas during the Silurian, from which the majority of eurypterid species have been described. The Silurian genus Eurypterus accounts for more than 90% of all known eurypterid specimens. Though the group continued to diversify during the subsequent Devonian period, the eurypterids were heavily affected by the Late Devonian extinction event. They declined in numbers and diversity until becoming extinct during the Permian–Triassic extinction event (or sometime shortly before) 251.9million years ago.
Although popularly called "sea scorpions", only the earliest eurypterids were marine; many later forms lived in brackish or fresh water, and they were not true scorpions. Some studies suggest that a dual respiratory system was present, which would have allowed for short periods of time in terrestrial environments. The name Eurypterida comes from the Ancient Greek words (), meaning 'broad' or 'wide', and (), meaning 'wing', referring to the pair of wide swimming appendages present in many members of the group.
The eurypterid order includes the largest known arthropods ever to have lived. The largest, Jaekelopterus, reached in length. Eurypterids were not uniformly large and most species were less than long; the smallest eurypterid, Alkenopterus, was only long. Eurypterid fossils have been recovered from every continent. A majority of fossils are from fossil sites in North America and Europe because the group lived primarily in the waters around and within the ancient supercontinent of Euramerica. Only a handful of eurypterid groups spread beyond the confines of Euramerica and a few genera, such as Adelophthalmus and Pterygotus, achieved a cosmopolitan distribution with fossils being found worldwide.
Description
Like all other arthropods, eurypterids possessed segmented bodies and jointed appendages (limbs) covered in a cuticle composed of proteins and chitin. As in other chelicerates, the body was divided into two tagmata (sections); the frontal prosoma (head) and posterior opisthosoma (abdomen). The prosoma was covered by a carapace (sometimes called the "prosomal shield") on which both compound eyes and the ocelli (simple eye-like sensory organs) were located.
The prosoma also bore six pairs of appendages which are usually referred to as appendage pairs I to VI. The first pair of appendages, the only pair placed before the mouth, is called the chelicerae (homologous to the fangs of spiders). They were equipped with small pincers used to manipulate food fragments and push them into the mouth. In one lineage, the Pterygotidae, the chelicerae were large and long, with strong, well-developed teeth on specialised chelae (claws). The subsequent pairs of appendages, numbers II to VI, possessed gnathobases (or "tooth-plates") on the coxae (limb segments) used for feeding. These appendages were generally walking legs that were cylindrical in shape and were covered in spines in some species. In most lineages, the limbs tended to get larger the farther back they were. In the Eurypterina suborder, the larger of the two eurypterid suborders, the sixth pair of appendages was also modified into a swimming paddle to aid in traversing aquatic environments.
The opisthosoma comprised 12 segments and the telson, the posteriormost division of the body, which in most species took the form of a blade-like shape. In some lineages, notably the Pterygotioidea, the Hibbertopteridae and the Mycteroptidae, the telson was flattened and may have been used as a rudder while swimming. Some genera within the superfamily Carcinosomatoidea, notably Eusarcana, had a telson similar to that of modern scorpions and may have been capable of using it to inject venom. The coxae of the sixth pair of appendages were overlaid by a plate that is referred to as the metastoma, originally derived from a complete exoskeleton segment. The opisthosoma itself can be divided either into a "mesosoma" (comprising segments 1 to 6) and "metasoma" (comprising segments 7 to 12) or into a "preabdomen" (generally comprising segments 1 to 7) and "postabdomen" (generally comprising segments 8 to 12).
The underside of the opisthosoma was covered in structures evolved from modified opisthosomal appendages. Throughout the opisthosoma, these structures formed plate-like structures termed ( in German). These created a branchial chamber (gill tract) between preceding and the ventral surface of the opisthosoma itself, which contained the respiratory organs. The second to sixth opisthosomal segments also contained oval or triangular organs that have been interpreted as organs that aid in respiration. These organs, termed or "gill tracts", would potentially have aided eurypterids to breathe air above water, while , similar to organs in modern horseshoe crabs, would cover the parts that serve for underwater respiration.
The appendages of opisthosomal segments 1 and 2 (the seventh and eighth segments overall) were fused into a structure termed the genital operculum, occupying most of the underside of the opisthosomal segment 2. Near the anterior margin of this structure, the genital appendage (also called the or the median abdominal appendage) protruded. This appendage, often preserved very prominently, has consistently been interpreted as part of the reproductive system and occurs in two recognized types, assumed to correspond to male and female.
Biology
Size
Eurypterids were highly variable in size, depending on factors such as lifestyle, living environment and taxonomic affinity. Sizes around are common in most eurypterid groups. The smallest eurypterid, Alkenopterus burglahrensis, measured just in length.
The largest eurypterid, and the largest known arthropod ever to have lived, is Jaekelopterus rhenaniae. A chelicera from the Emsian Klerf Formation of Willwerath, Germany measured in length, but is missing a quarter of its length, suggesting that the full chelicera would have been long. If the proportions between body length and chelicerae match those of its closest relatives, where the ratio between claw size and body length is relatively consistent, the specimen of Jaekelopterus that possessed the chelicera in question would have measured between , an average , in length. With the chelicerae extended, another meter (3.28 ft) would be added to this length. This estimate exceeds the maximum body size of all other known giant arthropods by almost half a meter (1.64 ft) even if the extended chelicerae are not included. Two other eurypterids have also been estimated to have reached lengths of 2.5 metres; Erettopterus grandis (closely related to Jaekelopterus) and Hibbertopterus wittebergensis, but E. grandis is very fragmentary and the H. wittenbergensis size estimate is based on trackway evidence, not fossil remains.
The family of Jaekelopterus, the Pterygotidae, is noted for several unusually large species. Both Acutiramus, whose largest member A. bohemicus measured , and Pterygotus, whose largest species P. grandidentatus measured , were gigantic. Several different contributing factors to the large size of the pterygotids have been suggested, including courtship behaviour, predation and competition over environmental resources.
Giant eurypterids were not limited to the family Pterygotidae. An isolated long fossil metastoma of the carcinosomatoid eurypterid Carcinosoma punctatum indicates the animal would have reached a length of in life, rivalling the pterygotids in size. Another giant was Pentecopterus decorahensis, a primitive carcinosomatoid, which is estimated to have reached lengths of .
Typical of large eurypterids is a lightweight build. Factors such as locomotion, energy costs in molting and respiration, as well as the actual physical properties of the exoskeleton, limit the size that arthropods can reach. A lightweight construction significantly decreases the influence of these factors. Pterygotids were particularly lightweight, with most fossilized large body segments preserving as thin and unmineralized. Lightweight adaptations are present in other giant paleozoic arthropods as well, such as the giant millipede Arthropleura, and are possibly vital for the evolution of giant size in arthropods.
In addition to the lightweight giant eurypterids, some deep-bodied forms in the family Hibbertopteridae were also very large. A carapace from the Carboniferous of Scotland referred to the species Hibbertoperus scouleri measures wide. As Hibbertopterus was very wide compared to its length, the animal in question could possibly have measured just short of in length. More robust than the pterygotids, this giant Hibbertopterus would possibly have rivalled the largest pterygotids in weight, if not surpassed them, and as such be among the heaviest arthropods.
Locomotion
The two eurypterid suborders, Eurypterina and Stylonurina, are distinguished primarily by the morphology of their final pair of appendages. In the Stylonurina, this appendage takes the form of a long and slender walking leg, while in the Eurypterina, the leg is modified and broadened into a swimming paddle. Other than the swimming paddle, the legs of many eurypterines were far too small to do much more than allow them to crawl across the sea floor. In contrast, a number of stylonurines had elongated and powerful legs that might have allowed them to walk on land (similar to modern crabs).
A fossil trackway was discovered in Carboniferous-aged fossil deposits of Scotland in 2005. It was attributed to the stylonurine eurypterid Hibbertopterus due to a matching size (the trackmaker was estimated to have been about long) and inferred leg anatomy. It is the largest terrestrial trackway—measuring long and averaging in width—made by an arthropod found thus far. It is the first record of land locomotion by a eurypterid. The trackway provides evidence that some eurypterids could survive in terrestrial environments, at least for short periods of time, and reveals information about the stylonurine gait. In Hibbertopterus, as in most eurypterids, the pairs of appendages are different in size (referred to as a heteropodous limb condition). These differently sized pairs would have moved in phase, and the short stride length indicates that Hibbertopterus crawled with an exceptionally slow speed, at least on land. The large telson was dragged along the ground and left a large central groove behind the animal. Slopes in the tracks at random intervals suggest that the motion was jerky. The gait of smaller stylonurines, such as Parastylonurus, was probably faster and more precise.
The functionality of the eurypterine swimming paddles varied from group to group. In the Eurypteroidea, the paddles were similar in shape to oars. The condition of the joints in their appendages ensured their paddles could only be moved in near-horizontal planes, not upwards or downwards. Some other groups, such as the Pterygotioidea, would not have possessed this condition and were probably able to swim faster. Most eurypterines are generally agreed to have utilized a rowing type of propulsion similar to that of crabs and water beetles. Larger individuals may have been capable of underwater flying (or subaqueous flight) in which the motion and shape of the paddles are enough to generate lift, similar to the swimming of sea turtles and sea lions. This type of movement has a relatively slower acceleration rate than the rowing type, especially since adults have proportionally smaller paddles than juveniles. However, since the larger sizes of adults mean a higher drag coefficient, using this type of propulsion is more energy-efficient.
Some eurypterines, such as Mixopterus (as inferred from attributed fossil trackways), were not necessarily good swimmers. It likely kept mostly to the bottom, using its swimming paddles for occasional bursts of movements vertically, with the fourth and fifth pairs of appendages positioned backwards to produce minor movement forwards. While walking, it probably used a gait like that of most modern insects. The weight of its long abdomen would have been balanced by two heavy and specialized frontal appendages, and the center of gravity might have been adjustable by raising and positioning the tail.
Preserved fossilized eurypterid trackways tend to be large and heteropodous and often have an associated telson drag mark along the mid-line (as with the Scottish Hibbertopterus track). Such trackways have been discovered on every continent except for South America. In some places where eurypterid fossil remains are otherwise rare, such as in South Africa and the rest of the former supercontinent Gondwana, the discoveries of trackways both predate and outnumber eurypterid body fossils. Eurypterid trackways have been referred to several ichnogenera, most notably Palmichnium (defined as a series of four tracks often with an associated drag mark in the mid-line), wherein the holotype of the ichnospecies P. kosinkiorum preserves the largest eurypterid footprints known to date with the found tracks each being about in diameter. Other eurypterid ichnogenera include Merostomichnites (though it is likely that many specimens actually represent trackways of crustaceans) and Arcuites (which preserves grooves made by the swimming appendages).
Respiration
In eurypterids, the respiratory organs were located on the ventral body wall (the underside of the opisthosoma). , evolved from opisthosomal appendages, covered the underside and created a gill chamber where the "gill tracts" were located. Depending on the species, the eurypterid gill tract was either triangular or oval in shape and was possibly raised into a cushion-like state. The surface of this gill tract bore several spinules (small spines), which resulted in an enlarged surface area. It was composed of spongy tissue due to many invaginations in the structure.
Though the is referred to as a "gill tract", it may not necessarily have functioned as actual gills. In other animals, gills are used for oxygen uptake from water and are outgrowths of the body wall. Despite eurypterids clearly being primarily aquatic animals that almost certainly evolved underwater (some eurypterids, such as the pterygotids, would even have been physically unable to walk on land), it is unlikely the "gill tract" contained functional gills when comparing the organ to gills in other invertebrates and even fish. Previous interpretations often identified the eurypterid "gills" as homologous with those of other groups (hence the terminology), with gas exchange occurring within the spongy tract and a pattern of branchio-cardiac and dendritic veins (as in related groups) carrying oxygenated blood into the body. The primary analogy used in previous studies has been horseshoe crabs, though their gill structure and that of eurypterids are remarkably different. In horseshoe crabs, the gills are more complex and composed of many lamellae (plates) which give a larger surface area used for gas exchange. In addition, the gill tract of eurypterids is proportionally much too small to support them if it is analogous to the gills of other groups. To be functional gills, they would have to have been highly efficient and would have required a highly efficient circulatory system. It is considered unlikely, however, that these factors would be enough to explain the large discrepancy between gill tract size and body size.
It has been suggested instead that the "gill tract" was an organ for breathing air, perhaps actually being a lung, plastron or a pseudotrachea. Plastrons are organs that some arthropods evolved secondarily to breathe air underwater. This is considered an unlikely explanation since eurypterids had evolved in water from the start and they would not have organs evolved from air-breathing organs present. In addition, plastrons are generally exposed on outer parts of the body while the eurypterid gill tract is located behind the . Instead, among arthropod respiratory organs, the eurypterid gill tracts most closely resemble the pseudotracheae found in modern isopods. These organs, called pseudotracheae, because of some resemblance to the tracheae (windpipes) of air-breathing organisms, are lung-like and present within the pleopods (back legs) of isopods. The structure of the pseudotracheae has been compared to the spongy structure of the eurypterid gill tracts. It is possible the two organs functioned in the same way.
Some researchers have suggested that eurypterids may have been adapted to an amphibious lifestyle, using the full gill tract structure as gills and the invaginations within it as pseudotrachea. This mode of life may not have been physiologically possible, however, since water pressure would have forced water into the invaginations leading to asphyxiation. Furthermore, most eurypterids would have been aquatic their entire lives. No matter how much time was spent on land, organs for respiration in underwater environments must have been present. True gills, expected to have been located within the branchial chamber within the , remain unknown in eurypterids.
Ontogeny
Like all arthropods, eurypterids matured and grew through static developmental stages referred to as instars. These instars were punctuated by periods during which eurypterids went through ecdysis (molting of the cuticle) after which they underwent rapid and immediate growth. Some arthropods, such as insects and many crustaceans, undergo extreme changes over the course of maturing. Chelicerates, including eurypterids, are in general considered to be direct developers, undergoing no extreme changes after hatching (though extra body segments and extra limbs may be gained over the course of ontogeny in some lineages, such as xiphosurans and sea spiders). Whether eurypterids were true direct developers (with hatchlings more or less being identical to adults) or hemianamorphic direct developers (with extra segments and limbs potentially being added during ontogeny) has been controversial in the past.
Hemianamorphic direct development has been observed in many arthropod groups, such as trilobites, megacheirans, basal crustaceans and basal myriapods. True direct development has on occasion been referred to as a trait unique to arachnids. There have been few studies on eurypterid ontogeny as there is a general lack of specimens in the fossil record that can confidently be stated to represent juveniles. It is possible that many eurypterid species thought to be distinct actually represent juvenile specimens of other species, with paleontologists rarely considering the influence of ontogeny when describing new species.
Studies on a well-preserved fossil assemblage of eurypterids from the Pragian-aged Beartooth Butte Formation in Cottonwood Canyon, Wyoming, composed of multiple specimens of various developmental stages of eurypterids Jaekelopterus and Strobilopterus, revealed that eurypterid ontogeny was more or less parallel and similar to that of extinct and extant xiphosurans, with the largest exception being that eurypterids hatched with a full set of appendages and opisthosomal segments. Eurypterids were thus not hemianamorphic direct developers, but true direct developers like modern arachnids.
The most frequently observed change occurring through ontogeny (except for some genera, such as Eurypterus, which appear to have been static) is the metastoma becoming proportionally less wide. This ontogenetic change has been observed in members of several superfamilies, such as the Eurypteroidea, the Pterygotioidea and the Moselopteroidea.
Feeding
No fossil gut contents from eurypterids are known, so direct evidence of their diet is lacking. The eurypterid biology is particularly suggestive of a carnivorous lifestyle. Not only were many large (in general, most predators tend to be larger than their prey), but they had stereoscopic vision (the ability to perceive depth). The legs of many eurypterids were covered in thin spines, used both for locomotion and the gathering of food. In some groups, these spiny appendages became heavily specialized. In some eurypterids in the Carcinosomatoidea, forward-facing appendages were large and possessed enormously elongated spines (as in Mixopterus and Megalograptus). In derived members of the Pterygotioidea, the appendages were completely without spines, but had specialized claws instead. Other eurypterids, lacking these specialized appendages, likely fed in a manner similar to modern horseshoe crabs, by grabbing and shredding food with their appendages before pushing it into their mouth using their chelicerae.
Fossils preserving digestive tracts have been reported from fossils of various eurypterids, among them Carcinosoma, Acutiramus and Eurypterus. Though a potential anal opening has been reported from the telson of a specimen of Buffalopterus, it is more likely that the anus was opened through the thin cuticle between the last segment before the telson and the telson itself, as in modern horseshoe crabs.
Eurypterid coprolites discovered in deposits of Ordovician age in Ohio containing fragments of a trilobite and eurypterid Megalograptus ohioensis in association with full specimens of the same eurypterid species have been suggested to represent evidence of cannibalism. Similar coprolites referred to the species Lanarkopterus dolichoschelus from the Ordovician of Ohio contain fragments of jawless fish and fragments of smaller specimens of Lanarkopterus itself.
Though apex predatory roles would have been limited to the very largest eurypterids, smaller eurypterids were likely formidable predators in their own right just like their larger relatives.
Reproductive biology
As in many other entirely extinct groups, understanding and researching the reproduction and sexual dimorphism of eurypterids is difficult, as they are only known from fossilized shells and carapaces. In some cases, there might not be enough apparent differences to separate the sexes based on morphology alone. Sometimes two sexes of the same species have been interpreted as two different species, as was the case with two species of Drepanopterus (D. bembycoides and D. lobatus).
The eurypterid prosoma is made up of the first six exoskeleton segments fused together into a larger structure. The seventh segment (thus the first opisthosomal segment) is referred to as the metastoma and the eighth segment (distinctly plate-like) is called the operculum and contains the genital aperature. The underside of this segment is occupied by the genital operculum, a structure originally evolved from ancestral seventh and eighth pair of appendages. In its center, as in modern horseshoe crabs, is a genital appendage. This appendage, an elongated rod with an internal duct, is found in two distinct morphs, generally referred to as "type A" and "type B". These genital appendages are often preserved prominently in fossils and have been the subject of various interpretations of eurypterid reproduction and sexual dimorphism.
Type A appendages are generally longer than those of type B. In some genera they are divided into different numbers of sections, such as in Eurypterus where the type A appendage is divided into three but the type B appendage into only two. Such division of the genital appendage is common in eurypterids, but the number is not universal; for instance, the appendages of both types in the family Pterygotidae are undivided. The type A appendage is also armed with two curved spines called (lit. 'fork' in Latin). The presence of in the type B appendage is also possible and the structure may represent the unfused tips of the appendages. Located between the dorsal and ventral surfaces of the associated with the type A appendages is a set of organs traditionally described as either "tubular organs" or "horn organs". These organs are most often interpreted as spermathecae (organs for storing sperm), though this function is yet to be proven conclusively. In arthropods, spermathecae are used to store the spermatophore received from males. This would imply that the type A appendage is the female morph and the type B appendage is the male. Further evidence for the type A appendages representing the female morph of genital appendages comes in their more complex construction (a general trend for female arthropod genitalia). It is possible that the greater length of the type A appendage means that it was used as an ovipositor (used to deposit eggs). The different types of genital appendages are not necessarily the only feature that distinguishes between the sexes of eurypterids. Depending on the genus and species in question, other features such as size, the amount of ornamentation, and the proportional width of the body can be the result of sexual dimorphism. In general, eurypterids with type B appendages (males) appear to have been proportionally wider than eurypterids with type A appendages (females) of the same genera.
The primary function of the long, assumed female, type A appendages was likely to take up spermatophore from the substrate into the reproductive tract rather than to serve as an ovipositor, as arthropod ovipositors are generally longer than eurypterid type A appendages. By rotating the sides of the operculum, it would have been possible to lower the appendage from the body. Due to the way different plates overlay at its location, the appendage would have been impossible to move without muscular contractions moving around the operculum. It would have been kept in place when not it use. The on the type A appendages may have aided in breaking open the spermatophore to release the free sperm inside for uptake. The "horn organs," possibly spermathecae, are thought to have been connected directly to the appendage via tracts, but these supposed tracts remain unpreserved in available fossil material.
Type B appendages, assumed male, would have produced, stored and perhaps shaped spermatophore in a heart-shaped structure on the dorsal surface of the appendage. A broad genital opening would have allowed large amounts of spermatophore to be released at once. The long associated with type B appendages, perhaps capable of being lowered like the type A appendage, could have been used to detect whether a substrate was suitable for spermatophore deposition.
Evolutionary history
Origins
Until 1882 no eurypterids were known from before the Silurian. Contemporary discoveries since the 1880s have expanded the knowledge of early eurypterids from the Ordovician period. The earliest eurypterids known today, the megalograptid Pentecopterus, date from the Darriwilian stage of the Middle Ordovician, 467.3 million years ago. There are also reports of even earlier fossil eurypterids in the Fezouata Biota of Late Tremadocian (Early Ordovician) age in Morocco, but these have yet to be thoroughly studied, and are likely to be peytoiid appendages.
Pentecopterus was a relatively derived eurypterid, part of the megalograptid family within the carcinosomatoid superfamily. Its derived position suggests that most eurypterid clades, at least within the eurypterine suborder, had already been established at this point during the Middle Ordovician. The earliest known stylonurine eurypterid, Brachyopterus, is also Middle Ordovician in age. The presence of members of both suborders indicates that primitive stem-eurypterids would have preceded them, though these are so far unknown in the fossil record. The presence of several eurypterid clades during the Middle Ordovician suggests that eurypterids either originated during the Early Ordovician and experienced a rapid and explosive radiation and diversification soon after the first forms evolved, or that the group originated much earlier, perhaps during the Cambrian period.
As such, the exact eurypterid time of origin remains unknown. Though fossils referred to as "primitive eurypterids" have occasionally been described from deposits of Cambrian or even Precambrian age, they are not recognized as eurypterids, and sometimes not even as related forms, today. Some animals previously seen as primitive eurypterids, such as the genus Strabops from the Cambrian of Missouri, are now classified as aglaspidids or strabopids. The aglaspidids, once seen as primitive chelicerates, are now seen as a group more closely related to trilobites.
The fossil record of Ordovician eurypterids is quite poor. The majority of eurypterids once reportedly known from the Ordovician have since proven to be misidentifications or pseudofossils. Today only 11 species can be confidently identified as representing Ordovician eurypterids. These taxa fall into two distinct ecological categories; large and active predators from the ancient continent of Laurentia, and demersal (living on the seafloor) and basal animals from the continents Avalonia and Gondwana. The Laurentian predators, classified in the family Megalograptidae (comprising the genera Echinognathus, Megalograptus and Pentecopterus), are likely to represent the first truly successful eurypterid group, experiencing a small radiation during the Late Ordovician.
Silurian
Eurypterids were most diverse and abundant between the Middle Silurian and the Early Devonian, with an absolute peak in diversity during the Pridoli epoch, 423 to 419.2 million years ago, of the very latest Silurian. This peak in diversity has been recognized since the early twentieth century; of the approximately 150 species of eurypterids known in 1916, more than half were from the Silurian and a third were from the Late Silurian alone.
Though stylonurine eurypterids generally remained rare and low in number, as had been the case during the preceding Ordovician, eurypterine eurypterids experienced a rapid rise in diversity and number. In most Silurian fossil beds, eurypterine eurypterids account for 90% of all eurypterids present. Though some were likely already present by the Late Ordovician (simply missing from the fossil record so far), a vast majority of eurypterid groups are first recorded in strata of Silurian age. These include both stylonurine groups such as the Stylonuroidea, Kokomopteroidea and Mycteropoidea as well as eurypterine groups such as the Pterygotioidea, Eurypteroidea and Waeringopteroidea.
The most successful eurypterid by far was the Middle to Late Silurian Eurypterus, a generalist, equally likely to have engaged in predation or scavenging. Thought to have hunted mainly small and soft-bodied invertebrates, such as worms, species of the genus (of which the most common is the type species, E. remipes) account for more than 90% (perhaps as many as 95%) of all known fossil eurypterid specimens. Despite their vast number, Eurypterus are only known from a relatively short temporal range, first appearing during the Late Llandovery epoch (around 432 million years ago) and being extinct by the end of the Pridoli epoch. Eurypterus was also restricted to the continent Euramerica (composed of the equatorial continents Avalonia, Baltica and Laurentia), which had been completely colonized by the genus during its merging and was unable to cross the vast expanses of ocean separating this continent from other parts of the world, such as the southern supercontinent Gondwana. As such, Eurypterus was limited geographically to the coastlines and shallow inland seas of Euramerica.
During the Late Silurian the pterygotid eurypterids, large and specialized forms with several new adaptations, such as large and flattened telsons capable of being used as rudders, and large and specialized chelicerae with enlarged pincers for handling (and potentially in some cases killing) prey appeared. Though the largest members of the family appeared in the Devonian, large two meter (6.5+ ft) pterygotids such as Acutiramus were already present during the Late Silurian. Their ecology ranged from generalized predatory behavior to ambush predation and some, such as Pterygotus itself, were active apex predators in Late Silurian marine ecosystems. The pterygotids were also evidently capable of crossing oceans, becoming one of only two eurypterid groups to achieve a cosmopolitan distribution.
Devonian
Though the eurypterids continued to be abundant and diversify during the Early Devonian (for instance leading to the evolution of the pterygotid Jaekelopterus, the largest of all arthropods), the group was one of many heavily affected by the Late Devonian extinction. The extinction event, only known to affect marine life (particularly trilobites, brachiopods and reef-building organisms) effectively crippled the abundance and diversity previously seen within the eurypterids.
A major decline in diversity had already begun during the Early Devonian and eurypterids were rare in marine environments by the Late Devonian. During the Frasnian stage four families went extinct, and the later Famennian saw an additional five families going extinct. As marine groups were the most affected, the eurypterids were primarily impacted within the eurypterine suborder. Only one group of stylonurines (the family Parastylonuridae) went extinct in the Early Devonian. Only two families of eurypterines survived into the Late Devonian at all (Adelophthalmidae and Waeringopteridae). The eurypterines experienced their most major declines in the Early Devonian, during which over 50% of their diversity was lost in just 10 million years. Stylonurines, on the other hand, persisted through the period with more or less consistent diversity and abundance but were affected during the Late Devonian, when many of the older groups were replaced by new forms in the families Mycteroptidae and Hibbertopteridae.
It is possible that the catastrophic extinction patterns seen in the eurypterine suborder were related to the emergence of more derived fish. Eurypterine decline began at the point when jawless fish first became more developed and coincides with the emergence of placoderms (armored fish) in both North America and Europe.
Stylonurines of the surviving hibbertopterid and mycteroptid families completely avoided competition with fish by evolving towards a new and distinct ecological niche. These families experienced a radiation and diversification through the Late Devonian and Early Carboniferous, the last ever radiation within the eurypterids, which gave rise to several new forms capable of "sweep-feeding" (raking through the substrate in search of prey).
Carboniferous and Permian
Only three eurypterid families—Adelophthalmidae, Hibbertopteridae and Mycteroptidae—survived the extinction event in its entirety. It was assumed that these were all freshwater animals, which would have rendered the eurypterids extinct in marine environments, and with marine eurypterid predators gone, sarcopterygians, such as the rhizodonts, were the new apex predators in marine environments. However, various recent findings raise doubts about this, and suggest that these eurypterids were euryhaline forms that lived in marginal marine environments, such as estuaries, deltas, lagoons, and coastal ponds. One argument is paleobiogeographical; pterygotoid distribution seems to require oceanic dispersal. A recent review of Adelophthalmoidea admitted that "There is much more of a marine influence in many of the sections yielding Adelophthalmus than has previously been acknowledged." Similarly, a study of the eurypterid Hibbertopterus from the Carboniferous of New Mexico concluded that the habitat of some eurypterids "may need to be re-evaluated". The sole surviving eurypterine family, Adelophthalmidae, was represented by only a single genus, Adelophthalmus. The hibbertopterids, mycteroptids and Adelophthalmus survived into the Permian.
Adelophthalmus became the most common of all late Paleozoic eurypterids, existing in greater number and diversity than surviving stylonurines, and diversified in the absence of other eurypterines. Out of the 33 species referred to Adelophthalmus, 23 (69%) are from the Carboniferous alone. The genus reached its peak diversity in the Late Carboniferous. Though Adelophthalmus had already been relatively widespread and represented around all major landmasses in the Late Devonian, the amalgamation of Pangaea into a global supercontinent over the course of the last two periods of the Paleozoic allowed Adelophthalmus to gain an almost worldwide distribution.
During the Late Carboniferous and Early Permian Adelophthalmus was widespread, living primarily in brackish and freshwater environments adjacent to coastal plains. These environments were maintained by favorable climate conditions. They did not persist as climate changes owing to Pangaea's formation altered depositional and vegetational patterns across the world. With their habitat gone, Adelophthalmus dwindled in number and had already gone extinct by the Leonardian stage of the Early Permian.
Mycteroptids and hibbertopterids continued to survive for some time, with one genus of each group known from Permian strata: Hastimima and Campylocephalus respectively. Hastimima went extinct during the Early Permian, as Adelophthalmus had, while Campylocephalus persisted longer. A massive incomplete carapace from Permian deposits in Russia represents the sole fossil remains of the species C. permianus, which might have reached in length, while originally considered to be Late Permian in age other sources suggest a Kungurian-Roadian age (~283-267 million years ago) for the specimen. A fragment of eurypterid cuticle, given the species name Woodwardopterus freemanorum is known from the terminal Permian of Australia, which represents the youngest known eurypterid. No eurypterids are known from fossil beds higher than the Permian. This indicates that the last eurypterids died either in the catastrophic extinction event at its end or at some point shortly before it. This extinction event, the Permian–Triassic extinction event, is the most devastating mass extinction recorded, and rendered many other successful Paleozoic groups, such as the trilobites, extinct.
History of study
The first known eurypterid specimen was discovered in the Silurian-aged rocks of New York, to this day one of the richest eurypterid fossil locations. Samuel L. Mitchill described the specimen, discovered near Westmoreland in Oneida county in 1818. He erroneously identified the fossil as an example of the fish Silurus, likely due to the strange, catfish-like appearance of the carapace. Seven years later, in 1825, James E. DeKay examined the fossil and recognized it as clearly belonging to an arthropod. He thought the fossil, which he named Eurypterus remipes, represented a crustacean of the order Branchiopoda, and suggested it might represent a missing link between the trilobites and more derived branchiopods. The name Eurypterus derives from Greek ( 'broad, wide') and ( 'wing').
In 1843, Hermann Burmeister published his view on trilobite taxonomy and how the group related to other organisms, living and extinct, in the work Die Organisation der Trilobiten aus ihren lebenden Verwandten entwickelt. He considered the trilobites to be crustaceans, as previous authors had, and classified them together with what he assumed to be their closest relatives, Eurypterus and the genus Cytherina, within a clade he named "Palaeadae". Within Palaeadae, Burmeister erected three families; the "Trilobitae" (composed of all trilobites), the "Cytherinidae" (composed only of Cytherina, an animal today seen as an ostracod) and the Eurypteridae (composed of Eurypterus, then including three species).
The fourth eurypterid genus to be described (following Hibbertopterus in 1836 and Campylocephalus in 1838, not identified as eurypterids until later), out of those still seen as taxonomically valid in modern times, was Pterygotus (), described by Louis Agassiz in 1839. Pterygotus was considerably larger in size than Eurypterus and when the first fossils were discovered by quarrymen in Scotland they were referred to as "Seraphims" by the workers. Agassiz first thought the fossils represented remains of fish, only recognizing their nature as arthropod remains five years later in 1844.
In 1849, Frederick M'Coy classified Pterygotus together with Eurypterus and Belinurus (a genus today seen as a xiphosuran) within Burmeister's Eurypteridae. M'Coy considered the Eurypteridae to be a group of crustaceans within the order Entomostraca, closely related to horseshoe crabs. A fourth genus, Slimonia, based on fossil remains previously assigned to a new species of Pterygotus, was referred to the Eurypteridae in 1856 by David Page.
Jan Nieszkowski's De Euryptero Remipede (1858) featured an extensive description of Eurypterus fischeri (now seen as synonymous with another species of Eurypterus, E. tetragonophthalmus), which, along with the monograph On the Genus Pterygotus by Thomas Henry Huxley and John William Salter, and an exhaustive description of the various eurypterids of New York in Volume 3 of the Palaeontology of New York (1859) by James Hall, contributed massively to the understanding of eurypterid diversity and biology. These publications were the first to fully describe the whole anatomy of eurypterids, recognizing the full number of prosomal appendages and the number of preabdominal and postabdominal segments. Both Nieszkowski and Hall recognized that the eurypterids were closely related to modern chelicerates, such as horseshoe crabs.
In 1865, Henry Woodward described the genus Stylonurus (named and figured, but not thoroughly described, by David Page in 1856) and raised the rank of the Eurypteridae to that of order, effectively creating the Eurypterida as the taxonomic unit it is seen as today. In the work Anatomy and Relations of the Eurypterida (1893), Malcolm Laurie added considerably to the knowledge and discussion of eurypterid anatomy and relations. He focused on how the eurypterids related to each other and to trilobites, crustaceans, scorpions, other arachnids and horseshoe crabs. The description of Eurypterus fischeri by Gerhard Holm in 1896 was so elaborate that the species became one of the most completely known of all extinct animals, so much so that the knowledge of E. fischeri was comparable with the knowledge of its modern relatives (such as the Atlantic horseshoe crab). The description also helped solidify the close relationship between the eurypterids and other chelicerates by showcasing numerous homologies between the two groups.
In 1912, John Mason Clarke and Rudolf Ruedemann published The Eurypterida of New York in which all eurypterid species thus far recovered from fossil deposits there were discussed. Clarke and Ruedemann created one of the first phylogenetic trees of eurypterids, dividing the order into two families; Eurypteridae (distinguished by smooth eyes and including Eurypterus, Anthraconectes, Stylonurus, Eusarcus, Dolichopterus, Onychopterus and Drepanopterus) and Pterygotidae (distinguished by faceted eyes and including Pterygotus, Erettopterus, Slimonia and Hughmilleria). Both families were considered to be descended from a common ancestor, Strabops. In line with earlier authors, Clarke and Ruedemann also supported a close relationship between the eurypterids and the horseshoe crabs (united under the class Merostomata) but also discussed alternative hypotheses such as a closer relation to arachnids.
Classification
Historically, a close relationship between eurypterids and xiphosurans (such as the modern Atlantic horseshoe crab) has been assumed by most researchers. Several homologies encourage this view, such as correlating segments of the appendages and the prosoma. Additionally, the presence of plate-like appendages bearing the "gill tracts" on appendages of the opisthosoma (the ) was cited early as an important homology. In the last few decades of the nineteenth century, further homologies were established, such as the similar structures of the compound eyes of Pterygotus and horseshoe crabs (seen as especially decisive as the eye of the horseshoe crab was seen as possessing an almost unique structure) and similarities in the ontogeny within both groups. These ontogenetical similarities were seen as most apparent when studying the nepionic stages (the developmental stage immediately following the embryonic stage) in both groups, during which both xiphosurans and eurypterids have a proportionally larger carapace than adults, are generally broader, possess a distinct ridge down the middle, have a lesser number of segments which lack differentiation and have an underdeveloped telson.
Due to these similarities, the xiphosurans and eurypterids have often been united under a single class or subclass called Merostomata (erected to house both groups by Henry Woodward in 1866). Though xiphosurans (like the eurypterids) were historically seen as crustaceans due to their respiratory system and their aquatic lifestyle, this hypothesis was discredited after numerous similarities were discovered between the horseshoe crabs and the arachnids. Some authors, such as John Sterling Kingsley in 1894, classified the Merostomata as a sister group to the Arachnida under the class "Acerata" within a subphylum "Branchiata". Others, such as Ray Lankester in 1909, went further and classified the Merostomata as a subclass within the Arachnida, raised to the rank of class.
In 1866, Ernst Haeckel classified the Merostomata (containing virtually only the Eurypterida) and Xiphosura within a group he named Gigantostraca within the crustaceans. Though Haeckel did not designate any taxonomic rank for this clade, it was interpreted as equivalent to the rank of subclass, such as the Malacostraca and Entomostraca, by later researchers such as John Sterling Kinsgsley. In subsequent research, Gigantostraca has been treated as synonymous with Merostomata (rarely) and Eurypterida itself (more commonly).
A phylogenetic analysis (the results presented in a cladogram below) conducted by James Lamsdell in 2013 on the relationships within the Xiphosura and the relations to other closely related groups (including the eurypterids, which were represented in the analysis by genera Eurypterus, Parastylonurus, Rhenopterus and Stoermeropterus) concluded that the Xiphosura, as presently understood, was paraphyletic (a group sharing a last common ancestor but not including all descendants of this ancestor) and thus not a valid phylogenetic group. Eurypterids were recovered as closely related to arachnids instead of xiphosurans, forming the group Sclerophorata within the clade Dekatriata (composed of sclerophorates and chasmataspidids). Lamsdell noted that it is possible that Dekatriata is synonymous with Sclerophorata as the reproductive system, the primary defining feature of sclerophorates, has not been thoroughly studied in chasmataspidids. Dekatriata is, in turn, part of the Prosomapoda, a group including the Xiphosurida (the only monophyletic xiphosuran group) and other stem-genera.
Internal relationships
The internal classification of eurypterids within the Eurypterida is based mainly on eleven established characters. These have been used throughout the history of eurypterid research to establish clades and genera. These characters include: the shape of the prosoma, the shape of the metastoma, the shape and position of the eyes, the types of prosomal appendages, the types of swimming leg paddles, the structure of the doublure (the fringe of the dorsal exoskeleton), the structure of the opithosoma, the structure of the genital appendages, the shape of the telson and the type of ornamentation present. It is worth noting that not all of these characters are of equal taxonomic importance. They are not applicable to all eurypterids either; stylonurine eurypterids lack swimming leg paddles entirely. Some characters, including the prosoma and metastoma shapes and the position and shapes of the eyes, are seen as important only for the distinction between different genera. Most superfamilies and families are defined based on the morphology of the appendages.
The most important character used in eurypterid taxonomy is the type of prosomal appendages as this character is used to define entire suborders. General leg anatomy can also be used to define superfamilies and families. Historically, the chelicerae were considered the most important appendages from a taxonomical standpoint since they only occurred in two general types: a eurypterid type with small and toothless pincers and a pterygotid type with large pincers and teeth. This distinction has historically been used to divide the Eurypterida into the two suborders Eurypterina (small chelicerae) and "Pterygotina" (large and powerful chelicerae). This classification scheme is not without problems. In Victor Tollerton's 1989 taxonomic revision of the Eurypterida, with suborders Eurypterina and Pterygotina recognized, several clades of eurypterids today recognized as stylonurines (including hibbertopterids and mycteroptids) were reclassified as non-eurypterids in the new separate order "Cyrtoctenida" on the grounds of perceived inconsistencies in the prosomal appendages.
Modern research favors a classification into suborders Eurypterina and Stylonurina instead, supported by phylogenetic analyses. In particular, pterygotid eurypterids share a number of homologies with derived eurypterine eurypterids such as the adelophthalmids, and are thus best classified as derived members of the same suborder. In the Stylonurina, the sixth pair of appendages is represented by long and slender walking legs and lack a modified spine (referred to as the podomere 7a). In most eurypterids in the Eurypterina, the sixth pair of appendages is broadened into swimming paddles and always has a podomere 7a. 75% of eurypterid species are eurypterines and they represent 99% of all fossil eurypterid specimens. Of all eurypterid clades, the Pterygotioidea is the most species-rich, with over 50 species. The second most species-rich clade is the Adelophthalmoidea, with over 40 species.
The cladogram presented below, covering all currently recognized eurypterid families, follows a 2007 study by O. Erik Tetlie. The stylonurine suborder follows a 2010 study by James Lamsdell, Simon J. Braddy and Tetlie. The superfamily "Megalograptoidea", recognized by Tetlie in 2007 and then placed between the Onychopterelloidea and Eurypteroidea, has been omitted as more recent studies suggest that the megalograptids were members of the superfamily Carcinosomatoidea. As such, the phylogeny of the Carcinosomatoidea follows a 2015 study by Lamsdell and colleagues.
| Biology and health sciences | Chelicerata | null |
765691 | https://en.wikipedia.org/wiki/Portable%20toilet | Portable toilet | A portable or mobile toilet (colloquial terms: thunderbox, porta-john, porta-potty or porta-loo) is any type of toilet that can be moved around, some by one person, some by mechanical equipment such as a truck and crane. Most types do not require any pre-existing services or infrastructure, such as sewerage, and are completely self-contained. The portable toilet is used in a variety of situations, for example in urban slums of developing countries, at festivals, for camping, on boats, on construction sites, and at film locations and large outdoor gatherings where there are no other facilities. Most portable toilets are unisex single units with privacy ensured by a simple lock on the door. Some portable toilets are small molded plastic or fiberglass portable rooms with a lockable door and a receptacle to catch the human excreta in a container.
A portable toilet is not connected to a hole in the ground (like a pit latrine), nor to a septic tank, nor is it plumbed into a municipal system leading to a sewage treatment plant. The chemical toilet is probably the most well-known type of portable toilet, but other types also exist, such as urine-diversion dehydration toilets, composting toilets, container-based toilets, bucket toilets, freezing toilets and incineration toilets. A bucket toilet is a very simple type of portable toilet.
Types
Chemical toilets
A chemical toilet collects human waste in a holding tank and uses chemicals to minimize the odors. These chemicals may either mask the odor or contain biocides that hinder odor-causing bacteria from multiplying, keeping the smell to a minimum. Chemical toilets include those on plane and trains (although many of these are now vacuum toilets), as well as much simpler ones.
Portable camping toilets
A simpler type of portable toilet may be used in travel trailers (caravans, camper vans) and on small boats. They are also referred to as "cassette toilet" or "camping toilet", or under brand names that have become generic trademarks. The Oxford English Dictionary lists "Porta Potti" ("with arbitrary respelling") as "A proprietary name for: a portable chemical toilet, as used by campers", and gives mostly American examples from 1968. The OED gives this proprietary name a second meaning, "a small prefabricated unit containing a toilet, designed for easy transportation and temporary installation esp. outdoors", which Wikipedia covers under chemical toilet.
The other name common in British English is "Elsan", which dates back to 1925. According to the Camping and Caravanning Club, "Today you will often see campsites refer to their Chemical Disposal Points as Elsan Disposal Points because of the history and popularity of the brand." The Canal and River Trust uses both brand names, in lieu of any unbranded term.
One colloquialism for these simple toilets is the "bucket and chuck it" system, although in fact they no longer resemble an open bucket (see bucket toilet). These are designed to be emptied into sanitary stations connected to the regular sewage system. These toilets are not to be confused with the types that are plumbed in to the vehicle and need to be pumped out at holding tank dump stations.
Urine-diversion dehydration toilets
Portable urine-diversion dehydration toilets are self-contained dry toilets sometimes referred to as "mobile" or "stand-alone" units. They are identifiable by their one-piece molded plastic shells or, in the case of DIY versions, simple plywood box construction. Most users of self-contained UDDTs rely upon a post-treatment process to ensure pathogen reduction. This post-treatment may consist of long-term storage or addition to an existing or purpose-built compost pile or some combination thereof. The necessity of a post-treatment step hinges upon the frequency and volume of use. For instances of infrequent or very modest seasonal use, a post-treatment phase might be deemed unnecessary due to the lower accumulation of waste, simplifying the overall disposal process.
Others
A commode chair (a chair enclosing a chamber pot) is a basic portable toilet that was used, for example, in 19th-century Europe.
History
The close stool, built as an article of furniture, is one of the earliest forms of portable toilet. They can still be seen in historic house museums such as Sir George-Étienne Cartier National Historic Site in Old Montreal, Canada. The velvet upholstered close stool used by William III is on display at Hampton Court Palace; see Groom of the Stool.
Early versions of the "Elsan chemical closet" ("closet" meaning a small room, see water closet, WC, and earth closet) were sold at Army & Navy Stores. Their use in World War II bomber aircraft is described at some length by the Bomber Command Museum of Canada; in brief, they were not popular with either the flying crew or the ground crew.
African-Americans living under Jim Crow laws (i.e. before the Civil Rights Act of 1964) faced severe challenges. Public toilets were segregated by race, and many restaurants and gas stations refused to serve black people, so some travellers carried a portable toilet in the trunk of their car.
Since 1974, Grand Canyon guides rafting on the Colorado River have used ammo boxes as portable toilets, typically with a removable toilet seat, according to the Museum of Northern Arizona in Flagstaff, Arizona.
Society and culture
A slang term, now dated or historic, is a "thunder-box" (Oxford English Dictionary: "a portable commode; by extension, any lavatory"). The term was used particularly in British India; travel writer Stephen McClarence called it "a crude sort of colonial lavatory". One features to comic effect in Evelyn Waugh's novel Men at Arms:
| Technology | Hydraulics and pneumatics | null |
765888 | https://en.wikipedia.org/wiki/Setter | Setter | The setter is a type of gundog used most often for hunting game such as quail, pheasant, and grouse.
In the UK, the four setter breeds, together with the pointers, usually form a subgroup within the gundog group as they share a common function. However, the setter breeds each have subtle differences in head, bone and substance.
The American and Canadian Kennel Clubs classify these breeds within the Sporting Group. Setters from show lines are usually considered to be heavier and larger than those from 'working' lines.
Function
A setter silently searches for game by scent; hunting is done systematically and methodically. When prey is encountered, the dog becomes motionless rather than chasing after the game. Setters get their name from their distinctive stance; a sort of crouch or "set" upon finding their quarry. Once the dog has indicated where the birds are by freezing on point, the birds are then flushed so the following guns can get a shot. In earlier times before guns were used, a net would be used to trap the birds.
The scent of game birds is airborne, so to sense it, the setter carries its head high and should never follow foot scent. Most setters are born with a natural proclivity to hunting. Dogs which show excitement and interest in birds are described as being "birdy", and trainers look for puppies that show this particular trait. Training is usually done with quail as a first choice or domesticated pigeons.
Attributes
This group of dogs combines beauty, brains and bird sense; the early setter breeds are believed to have been developed as far back as the 15th century in the UK. The ancestors of modern setters probably originated in Spain and were bred from spaniel stock. Later, these dogs were exported to France and England where the breeds were developed into today's varieties.
They are fast, stylish game-finding dogs with a unique history and evolution for the single purpose of finding game birds. Writing in 1576 Dr Johannes Caius states "There is also at this date among us a new kind of dogge brought out of Fraunce, and they bee speckled all over with white and black, which mingled colours incline to a marble blewe". Argue speculates this may be a description of the blue belton colour found in English setters.
Early shows and field trials
The first official dog show held in the UK was at Newcastle-on-Tyne in June 1859 and entry was restricted to setters and pointers. There were 36 setters and 23 pointers entered. The show was organised by John Shorthose and William Pape. Mr Jobling's Black and Tan Setter, Dandy, won the first prize for setters. The class for pointers was judged by Mr Jobling who awarded the prize to a pointer owned by a Mr Brailsford, who helped judge the setters. This raised some criticism.
The prize awarded to each winner was a double barrelled gun worth around £15 to £20.
There was uncertainty as to how setters would be classified at early shows. Three classes were usually scheduled in 1862 dividing setters into three categories: English, Black/Tan and Irish. These became official breed classifications when The Kennel Club was founded in 1873.
During 1806 in the UK there was a sale of setters. A black setter bitch called Peg was sold for 41 guineas while the price for setter dogs called Punch, Brush, Bob, Bell, Bounce and Sam varied from 17 to 32 guineas each. No colours were specified for the dogs.
The first recorded field trial in the UK was held in April 1865 on the estate of Samuel Whitbread (MP) at Southill, Bedfordshire. It was only open to setters and pointers. All the setters entered were black and tans (Gordons).
In 1879, the Western Hemisphere held its first recorded show. This was also restricted to setters and pointers. It was four years after this that the first American Field trial was held.
More recent competitions
The Kennel Club has four champion titles available to be achieved by setters competing in the UK. These are:
Show Champion (Sh Ch), awarded to dogs who have won three Challenge Certificates (CCs) under three different judges with at least one CC won after 12 months of age
Champion (Ch), gained by dogs which have won a Sh Ch title plus a field trial award, Diploma of Merit or a Show Gundog Working Certificate
Field Trial Champion (Ft Ch), dogs which won a pointer or setter open stake or two first prizes at two different stakes under two different A Panel judges; there must be no less than 16 runners entered.
Dual Champion, the highest award available to setters, a dog which has achieved the titles of Show Champion and Field Trial Champion.
Challenge Certificates were first introduced by the Kennel Club in December 1900. Prior to 1958, a gundog could not claim to be a champion no matter how many CCs it won until it had gained a qualifier in the field as well. When the Show Champion title was instigated in 1958, it was agreed this could be applied retroactively.
The first ever gundog to attain the title was an English Setter. Since the Second World War only two dogs have achieved Dual Champion status in the UK. The first was a Pointer and the second was a Gordon Setter, the only Gordon to ever achieve this accolade.
Setters have been Best in Show at Crufts seven times. The award was secured by Irish Setters in 1981, 1993, 1995 and 1999. English Setters were best in show in 1964, 1977 and 1988. The Irish Setter Best in Show of 1981 was already a Field Trial Champion, proving that she had brains as well as beauty.
At the Westminster show in America an English Setter won the Best in Show title in 1938. He was only 11 months old and at his very first show. This was before entry to the show was restricted to Champions in 1992. He is the only setter to achieve Best in Show at Westminster.
Registrations
In January 2006, the Kennel Club identified a number of British native breeds registering 300 or fewer puppies each year; it labelled these breeds as 'Vulnerable Native Breeds'. Initially, the list included the Irish Red and White Setter and the Gordon Setter but Gordons were re-classified as viable in January 2007 after consultation with the breed clubs. English Setters were added for the first time in 2012.
To give an indication of how the UK registrations have changed, English Setter registrations were 568 in 2002, Gordons were 250 and Irish Red and Whites were 99, but Irish registrations totalled 1,225. However, by 2011 the figures for English Setters dropped to 234 puppy registrations; Gordons had a slight increase to 306; Irish Red and Whites had a slight drop to 83; and Irish decreased to 869 puppy registrations.
In contrast, in a comment about registration figures and popular breed rankings, the American Kennel Club stated that 2011 was the "year of the setters, with all four making big jumps over the past year". English Setters had ranked at 101 in 2010 but moved up to 87 in 2011; Irish had shifted from ranking at 77 in 2010 to 70 in 2011; Irish Red and Whites went up three places having been 150 in 2010 and 147 in 2011; and Gordons improved its 2010 position of 98 to be ranked at 94.
On January 1, 2009, the Irish Red and White Setter became eligible for American Kennel Club registration and was thereafter able to compete in the Sporting Group at its shows.
Coat
Most setter breeds have long smooth, silky coats that require maintenance. While Gordon, Irish and English Setters usually undergo some trimming for presentation in the show ring, Irish Red and White Setters do not require as much work, because they have lighter coats.
Temperament
Setters have a tendency to be happy, playful dogs and are usually very friendly both to people and other dogs. They have a great deal of energy and require daily exercise.
Setter breeds
The breeds making up this subgroup are:
English Setter
Gordon Setter
Irish Setter (a.k.a. Irish Red Setter)
Irish Red and White Setter
| Biology and health sciences | Dogs | Animals |
765968 | https://en.wikipedia.org/wiki/Water%20dog | Water dog | A water dog is a type of gundog bred to flush and retrieve game from water. Water dogs are considered the progenitors of most modern retriever dog breeds.
Description
Water dogs are usually medium-sized, active dogs; their most distinctive feature are their tight waterproof coats and their strong desire to swim. Traditionally many long haired water dogs breeds have their coats clipped with a bare midriff and hindquarters to assist in swimming by reducing drag, whilst retaining a long coat around their torso to prevent thermal shock when jumping into freezing water. This classic clip is seen to this day in dog shows with some breeds such as Poodles retaining a variation of this clip. Paintings by artists like Francisco Goya and Albrecht Durer in the 15th century depict early Standard Poodles, which were used for retrieving bird games from water.
History
Water dogs are an ancient type of dog, used by ancient mariners to retrieve objects lost overboard as well as to pass messages between boats, they were known to the Romans who called them "lion dogs" after the distinctive clipping of their coats. Water dogs continued to be seen aboard ships until modern times. They were known to be prevalent in the Spanish Armada and could still be seen through to the 20th century aboard fishing vessels, particularly in the waters around Spain and Portugal.
In addition to their roles as ships' dogs, in Medieval Europe water dogs were used quite widely by hunters to retrieve ducks from water that had been shot with a bow and arrows, as well as the arrows and bolts that had missed their mark.
With the invention of modern firearms the need for dogs to retrieve arrows disappeared and many of the European water dogs breeds became predominantly ornamental, such as the poodle, whilst others adapted into modern gundogs, like the Wetterhoun. It is believed that water dogs were used, along with other types, in the breeding of most modern retriever breeds.
List of breeds
This is a list of dog breeds usually considered water dogs.
† Extinct breed‡ Often considered a versatile gundog
| Biology and health sciences | Dogs | Animals |
765972 | https://en.wikipedia.org/wiki/SRGB | SRGB | sRGB is a standard numerical encoding of colors, based on the RGB (red, green, blue) color space, for use on monitors, printers, and the World Wide Web. It was initially proposed by HP and Microsoft in 1996 and became an official standard of the International Electrotechnical Commission (IEC) as IEC 61966-2-1:1999. It is the current defined standard colorspace for the web, and it is usually the assumed colorspace for images that are neither tagged for a colorspace nor have an embedded color profile.
The sRGB standard uses the same color primaries and white point as ITU-R BT.709 standard for HDTV, but a different transfer function (or gamma) compatible with the era's CRT displays, and assumes a viewing environment closer to typical home and office viewing conditions.
The sRGB color space is also the basis of the sYCC color encoding, which is a remapping of the R, G, and B components of sRGB to a luminance (brightness) value and two chroma channels similar to those of the CIE YCbCr encoding.
sRGB definition
Gamut
The sRGB standard defines the chromaticities of the red, green, and blue primaries, the colors where one of the three channels is nonzero and the other two are zero. The gamut of chromaticities that can be represented in sRGB is the color triangle defined by these primaries, which are set such that the range of colors inside the triangle is well within the range of colors visible to a human with normal trichromatic vision. As with any RGB color space, for non-negative values of R, G, and B it is not possible to represent colors outside this triangle.
The primaries come from HDTV (ITU-R BT.709), which are somewhat different from those for older color TV systems (ITU-R BT.601). These values were chosen to reflect the approximate color of consumer CRT phosphors at the time of its design. Since flat-panel displays at the time were generally designed to emulate CRT characteristics, the values also reflected prevailing practice for other display devices as well.
Transfer function ("gamma")
The sRGB standard specifies a non-linear encoding of physical brightness values (proportional to luminous power emitted per unit of area) into the integer R, G, and B values that are to be stored in computer memory or image files. This transfer function commonly called gamma encoding, is the combination of a linear function at low brightness values and a displaced power law for the rest of the range.
Specifically, let be the encoded R, G, or B value, assumed to be an integer ranging from 0 (meaning no light) to some maximum (meaning the maximum displayable intensity of that channel). Typically is 255 when as an 8-bit integer, or generally for an -bit integer. The physical intensity represented by is defined as , where the decoding function is defined as
where , , , and . The result is 0 for no light, and 1 for the maximum intensity. The value is sometimes called the "linear value" or "linear-light value" corresponding to the encoded sample .
Conversely, given a value between 0 and 1 that is proportional to the physical R, G, or B intensity to be displayed, the encoded integer intensity will be , where encoding function is defined as
where , and , , and are the same as in the decoding function .
These sRGB encoding and decoding functions and above are similar in form to those of BT.709, but the values are noticeably different. Because of the rounding of the parameters, they have small discontinuities at the transition between the linear and non-linear part, on the order of 10−8 — which are too small to matter in most practical situations.
For medium and large values of the argument , the decoding function is actually quite close to a pure power law with exponent (gamma) of 2.2. However, for low values (around and below the breakpoint ) the difference is perceptible.
If needed, the encoding and decoding functions , can be used for arguments greater than 1. They are also extended to negative values by the identities , .
In practice, there is still debate and confusion about the formulas used for encoding and decoding image colors from or into "sRGB" files. Also, programmers may choose use the "pure" power law instead of the correct piecewise formulas above, for the sake of code simplicity or efficiency, but that would introduce some issues near black. Also, some operating systems and programs may directly send sRGB encoded images to monitors which actually have a different effective decoding function, such as pure power law with 2.2 exponent, this is further complicated by the standard saying that reference CRT display should use 2.2 gamma.
Correspondence to CIE XYZ stimulus
The sRGB standard specifies also the colors and relative intensities of the three primaries R, G, and B, by defining the mapping between these values (in linear brightness scale, before the gamma encoding) and the CIE XYZ perceptual color coordinates. This mapping is the same specified by the BT.709 standard; in matrix notation,
These coefficients should be considered exact and assume the 2° standard colorimetric observer for CIE XYZ. In particular, the second row of this matrix specifies the computation of the BT.709-2 luma (brightness) value from the linear R, G, and B values. (BT.709-1 had a typo in these coefficients.)
The inverse conversion, from from CIE XYZ to (linear) sRGB, can be obtained by inverting the matrix above to a suitable numerical accuracy. The 1999 standard provides the matrix
which is not the exact inverse of the sRGB to XYZ transformation, but was expected to be accurate enough for 8-bit encoded samples (with ).
The 2003 amended version of the sRGB standard points out that higher accuracy is needed when the samples are encoded with more than 8 bits. It provides the matrix
.
which is claimed to be sufficiently accurate for samples with bits ().
For these formulas, the X, Y, and Z values must be scaled so that the Y of D65 ("white") is 1.0 (X = 0.9505, Y = 1.0000, Z = 1.0890). This is usually true but some color spaces use 100 or other values (such as in CIELAB, when using specified white points).
Viewing environment
The sRGB specification assumes a dimly lit encoding (creation) environment with an ambient correlated color temperature (CCT) of 5003 K:
The assumed ambient CCT differs from that of the BT.709 standard illuminant (D65), which is still retained for the screen white point. Using D50 for both would have made the white point of most photographic paper appear excessively blue. The other parameters, such as the luminance level, are representative of a typical CRT monitor.
For optimal results, the ICC recommends using the encoding viewing environment (i.e., dim, diffuse lighting) rather than the less-stringent typical viewing environment.
The bg-sRGB space
Annex G of the 2003 amendment of the sRGB standard describes an alternative encoding of color values, called bg-sRGB, that is recommended when the number of bits per channel is 10 or more.
In this alternative encoding, an integer encoded sample value and the corresponding linear R, G, or B value are related by the formulas
and
where the and functions are the same as in the sRGB space. The standard specifies and for , and generally
for .
In the bg-sRGB color space, the null intensity is encoded as the integer , and the maximum displayable intensity is encoded as the integer . This encoding is useful for color space manipulations (like the conversion from sYCC) that can generate values that are negative or greater than 1.
History
The non-linear encoding of physical data samples is a common digital signal processing technique that aims to make better use of the bits available for the encoded signal, taking into account the non-linear way human senses perceive physical stimuli. Using smaller increments for smaller signals reduces the quantization artifacts.
This principle was incorporated into the digital-to-analog converters and the analog circuitry of early computer monitors, resulting in an effective decoding function (the mapping from digital sample values to the displayed intensity) which was roughly a power law with an exponent between 2 and 3. The exponent was commonly denoted with the letter , hence the common name "gamma correction" (or similar) for this function. This mapping initially varied according to CRT manufacturers, but was normalized in 1993 for use in HDTV systems, as the ITU BT.709 standard The BT.709 standard specified a decoding function with a linear section near zero, transitioning to a shifted power law with exponent 1/0.45 ≈ 2.2222....
The sRGB encoding was created a few years later by Hewlett-Packard and Microsoft. It was meant to describe the decoding function of most CRT computer monitors used with Windows operating systems at the time, which was still different from that assumed by BT.709. The first draft of the standard was published in 1996. A fourth draft, still incomplete, is available online. Like the BT.709, the sRGB decoding function was defined as a linear section near zero that transitions to a shifted power law
Justification for the formulas
In theory, the parameters of the encoding and decoding functions should be chosen so that the transition from the linear section to the power law section is continuous (without a sudden step) and smooth (without a sudden change of slope).
To derive the decoding function, one considers that the general formula for a linear function, whose graph is a straight line that passes through , is , and a shifted power law curve that passes through is
To obtain a seamless transition between the two functions when has a value , we must have
To avoid a sudden change of slope where the two segments meet, the derivatives must be equal at this point:
Solving the two equations for and we get
Parameter values
The first draft of the sRGB standard initially set the parameters at and so that the resulting decoding function closely resembled a pure power law with exponent (gamma) 2.2, assimed to be typical of computer monitors at the time. This choice implies a breakpoint and a linear coefficient . These values, rounded to and are still incorrectly given in some publications. However, the value of was rounded to already in the sRGB draft standard, resulting in a small discontinuity in the curve.
The first official version of the standard was defined and published by the IEC in 1999. In this version, the rounded value of was retained, but the breakpoint was redefined as to make the curve approximately continuous. With these values, there is still a discontinuity in the slope, from just below the intersection to just above it. The final standard also corrected some small rounding errors present in the draft.
The 1999 IEC standard was amended in 2003. The sRGB to CIE XYZ matrix was retained, but the inverse transformation above was replaced by a more accurate version, with seven decimal fraction digits. The amended standard also included the definition of the sYCC encoding, using brightness (Y) and two chroma coordinates (CC) instead of R, G, and B coordinates, and a bg-sRGB encoding for 10-bit components which allows some component values outside the 0-1 range.
Usage
Due to the standardization of sRGB on the Internet, on computers, and on printers, many low- to medium-end consumer digital cameras and scanners use sRGB as the default (or only available) working color space. However, consumer-level CCDs are typically uncalibrated, meaning that even though the image or device is being labeled as "sRGB", one cannot assume that the encoded values or the colors of displayed images are accurate as specified by the standard. Recently that changed and smartphones and TVs have very good calibration out of the box.
If the color space of an image is unknown and the R, G, and B samples are encoded with 8 bits each, the sRGB encoding usually the assumed default.
As the sRGB gamut mostly meets or exceeds the gamut of a low-end inkjet printer, an sRGB image is often regarded as satisfactory for home printing. The sRGB color space is sometimes avoided by high-end print publishing professionals because its color gamut is not big enough, especially in the blue-green colors, to include all the colors that can be reproduced in CMYK printing. Images intended for professional printing via a fully color-managed workflow (e.g. prepress output) sometimes use another color space such as Adobe RGB (1998), which accommodates a wider gamut and CMYK color space like Fogra39. Such CMYK and sRGB images used on the Internet may be converted to sRGB using color management tools that are usually included with software that works in these other color spaces.
Programming interface support
The two dominant programming interfaces for 3D graphics, OpenGL and Direct3D, have both incorporated support for the sRGB gamma curve. OpenGL supports textures with sRGB gamma encoded color components (first introduced with EXT_texture_sRGB extension, added to the core in OpenGL 2.1) and rendering into sRGB gamma encoded framebuffers (first introduced with EXT_framebuffer_sRGB extension, added to the core in OpenGL 3.0). Correct mipmapping and interpolation of sRGB gamma textures has direct hardware support in texturing units of most modern GPUs (for example nVidia GeForce 8 performs conversion from 8-bit texture to linear values before interpolating those values), and does not have any performance penalty.
ICC profiles
A lookup table may be used to efficiently convert sRGB to other color spaces. The International Color Consortium (ICC) has published color profiles for this purpose, which are widely used. There are several variants, including ICCmax, version 4, and version 2.
Version 4 is generally recommended, but version 2 is still commonly used and is the most compatible with other software including browsers. However, inconsistencies have been pointed out between those ICC profiles and the IEC sRGB standard. In particular, version 2 of the ICC profile specification does not implement the piecewise parametric curve encoding ("para") as specified by the IEC sRGB standard, and has to implement the linear spline using 1DLUT. What is worse in some cases "simplified sRGB" is used which is just 2.2 gamma.
The sYCC color space
Amendment 1 to IEC 61966-2-1:1999, approved in 2003, includes the definition of a Y′Cb′Cr′ color representation called sYCC. Although the RGB color primaries are based on BT.709, the equations for transformation from sRGB to sYCC and vice versa are based on BT.601. The sYCC standard specifies 8 bits for the encoded components, and the matrices result in a range of approximately 0–1 for Y; -0.5–0.5 for C.
As this conversion can result in sRGB values outside the range 0–1, the amendment describes how to apply the gamma correction to negative values, by applying when is negative (and is the sRGB↔linear functions described above). This is also used by scRGB.
| Physical sciences | Basics | Physics |
765992 | https://en.wikipedia.org/wiki/Adobe%20RGB%20color%20space | Adobe RGB color space | The Adobe RGB (1998) color space or opRGB is a color space developed by Adobe Inc. in 1998. It was designed to encompass most of the colors achievable on CMYK color printers, but by using RGB primary colors on a device such as a computer display. The Adobe RGB (1998) color space encompasses roughly 30% of the visible colors specified by the CIELAB color space – improving upon the gamut of the sRGB color space, primarily in cyan-green hues. It was subsequently standardized by the IEC as IEC 61966-2-5:1999 with a name opRGB (optional RGB color space) and is used in HDMI.
Historical background
Beginning in 1997, Adobe Systems was looking into creating ICC profiles that its consumers could use in conjunction with Photoshop's new color management features. Since not many applications at the time had any ICC color management, most operating systems did not ship with useful profiles.
Lead developer of Photoshop, Thomas Knoll decided to build an ICC profile around specifications he found in the documentation for the SMPTE 240M standard, the precursor to Rec. 709 (but not in primaries: 240M also defined EOTF and thus was display referred, sRGB was created by connecting BT.470 PAL and SMPTE C). SMPTE 240M's gamut is wider than that of the BT.709 gamut and the same as BT.470 NTSC (System B, G). However, with the release of Photoshop 5.0 nearing, Adobe made the decision to include the profile within the software.
Although users loved the wider range of reproducible colors, those familiar with the SMPTE 240M specifications contacted Adobe, informing the company that it had copied the values that described idealized primaries, not actual standard ones (in a special annex to the standard). The real values were much closer to sRGB's, which avid Photoshop consumers did not enjoy as a working environment. To make matters worse, an engineer had made an error when copying the red primary chromaticity coordinates, resulting in an even more inaccurate representation of the SMPTE standard. On the other hand red and blue primary are the same as in PAL and green is the same as in NTSC 1953 (blue is the same as in BT.709 and sRGB).
Adobe tried numerous tactics to correct the profile, such as correcting the red primary and changing the white point to match that of the CIE Standard Illuminant D50 (though that will also change the primaries and is thus pointless), yet all of the adjustments made CMYK conversion worse than before. In the end, Adobe decided to keep the "incorrect" profile, but changed the name to Adobe RGB (1998) in order to avoid a trademark search or infringement.
Specifications
Reference viewing conditions
In Adobe RGB (1998), colors are specified as [R,G,B] triplets, where each of the R, G, and B components have values ranging between 0 and 1. When displayed on a monitor, the exact chromaticities of the reference white point [1,1,1], the reference black point [0,0,0], and the primaries ([1,0,0], [0,1,0], and [0,0,1]) are specified. To meet the color appearance requirements of the color space, the luminance of the monitor must be 160.00 cd/m2 at the white point, and 0.5557 cd/m2 at the black point, which implies a contrast ratio of 287.9. Moreover, the black point shall have the same chromaticity as the white point, yet with a luminance equal to 0.34731% of the white point luminance. The ambient illumination level at the monitor faceplate when the monitor is turned off must be 32 lx.
As with sRGB, the RGB component values in Adobe RGB (1998) are not proportional to the luminances. Rather, a gamma of approximately 2.2 is assumed, without the linear segment near zero that is present in sRGB. The precise gamma value is 563/256, or 2.19921875. In coverage of the CIE 1931 color space the Adobe RGB (1998) color space covers 52.1%.
The chromaticities of the primary colors and the white point, both of which correspond to the CIE Standard Illuminant D65, are as follows:
The corresponding absolute XYZ tristimulus values for the reference display white and black points are as follows:
Normalized XYZ tristimulus values can be obtained from absolute luminance XaYaZa tristimulus values as follows:
where XKYKZK and XWYWZW are reference display black and white points in the table above.
The conversion between normalized XYZ to and from Adobe RGB tristimulus values can be done as follows:
As was later defined in the IEC standard opYCC uses BT.601 matrix for conversion to YCbCr, that can be full range matrix and limited range matrix. Display can signal YCC quantization range support and sink can send either one.
ICC PCS color image encoding
An image in the ICC Profile Connection Space (PCS) is encoded in 24-bit Adobe RGB (1998) color image encoding. Through the application of the 3x3 matrix below (derived from the inversion of the color space chromaticity coordinates and a chromatic adaptation to CIE Standard Illuminant D50 using the Bradford transformation matrix), the input image's normalized XYZ tristimulus values are transformed into RGB tristimulus values. The component values would be clipped to the range [0, 1].
The RGB tristimulus values are then converted to Adobe RGB R'G'B''' component values through the use of the following component transfer functions:
The resulting component values would be then represented in floating point or integer encodings. If it is necessary to encode values from the PCS back to the input device space, the following matrix can be implemented:
Comparison to sRGB
Gamut
sRGB is an RGB color space proposed by HP and Microsoft in 1996 to approximate the color gamut of the (then) most common computer display devices (CRTs). Since sRGB serves as a "best guess" metric for how another person's monitor produces color, it has become the standard color space for displaying images on the Internet. sRGB's color gamut encompasses just 35% of the visible colors specified by CIE, whereas Adobe RGB (1998) encompasses slightly more than 50% of all visible colors. Adobe RGB (1998) extends into richer cyans and greens than does sRGB – for all levels of luminance. The two gamuts are often compared in mid-tone values (~50% luminance), but clear differences are evident in shadows (~25% luminance) and highlights (~75% luminance) as well. In fact, Adobe RGB (1998) expands its advantages to areas of intense orange, yellow, and magenta regions.
Although there is a significant difference between gamut ranges in the CIE xy chromaticity diagram, if the coordinates were to be transformed to fit on the CIE u′v′ chromaticity diagram, which illustrates the eye's perceived variance in hue more closely, the difference in the green region is far less exaggerated. Also, although Adobe RGB (1998) can theoretically'' represent a wider gamut of colors, the color space requires special software and a complex workflow in order to utilize its full range. Otherwise, the produced colors would be squeezed into a smaller range (making them appear duller) in order to match sRGB's more widely used gamut.
Bit depth distribution
Although the Adobe RGB (1998) working space clearly provides more colors to utilize, another factor to consider when choosing between color spaces is how each space influences the distribution of the image's bit depth. Color spaces with larger gamuts "stretch" the bits over a broader region of colors, whereas smaller gamuts concentrate these bits within a narrow region.
A similar, yet not as dramatic concentration of bit depth occurs with Adobe RGB (1998) versus sRGB, except in three dimensions rather than one. The Adobe RGB (1998) color space occupies roughly 40% more volume than the sRGB color space, which concludes that one would only be exploiting 70% of the available bit depth if the colors in Adobe RGB (1998) are unnecessary. On the contrary, one may have plenty of "spare" bits if using a 16-bit image, thus negating any reduction due to the choice of working space.
| Physical sciences | Basics | Physics |
766619 | https://en.wikipedia.org/wiki/Cyclonic%20separation | Cyclonic separation | Cyclonic separation is a method of removing particulates from an air, gas or liquid stream, without the use of filters, through vortex separation. When removing particulate matter from liquid, a hydrocyclone is used; while from gas, a gas cyclone is used. Rotational effects and gravity are used to separate mixtures of solids and fluids. The method can also be used to separate fine droplets of liquid from a gaseous stream.
Operation
A high-speed rotating (air)flow is established within a cylindrical or conical container called a cyclone. Air flows in a helical pattern, beginning at the top (wide end) of the cyclone and ending at the bottom (narrow) end before exiting the cyclone in a straight stream through the center of the cyclone and out the top. Larger (denser) particles in the rotating stream have too much inertia to follow the tight curve of the stream, and thus strike the outside wall, then fall to the bottom of the cyclone where they can be removed. In a conical system, as the rotating flow moves towards the narrow end of the cyclone, the rotational radius of the stream is reduced, thus separating smaller and smaller particles. The cyclone geometry, together with volumetric flow rate, defines the cut point of the cyclone. This is the size of particle that will be removed from the stream with a 50% efficiency. Particles larger than the cut point will be removed with a greater efficiency, and smaller particles with a lower efficiency as they separate with more difficulty or can be subject to re-entrainment when the air vortex reverses direction to move in direction of the outlet.
An alternative cyclone design uses a secondary air flow within the cyclone to keep the collected particles from striking the walls, to protect them from abrasion. The primary air flow containing the particulates enters from the bottom of the cyclone and is forced into spiral rotation by stationary spinner vanes. The secondary air flow enters from the top of the cyclone and moves downward toward the bottom, intercepting the particulate from the primary air. The secondary air flow also allows the collector to optionally be mounted horizontally, because it pushes the particulate toward the collection area, and does not rely solely on gravity to perform this function.
Uses
Cyclone separators are found in all types of power and industrial applications, including pulp and paper plants, cement plants, steel mills, petroleum coke plants, metallurgical plants, saw mills and other kinds of facilities that process dust.
Large scale cyclones are used in sawmills to remove sawdust from extracted air. Cyclones are also used in oil refineries to separate oils and gases, and in the cement industry as components of kiln preheaters. Cyclones are increasingly used in the household, as the core technology in bagless types of portable vacuum cleaners and central vacuum cleaners. Cyclones are also used in industrial and professional kitchen ventilation for separating the grease from the exhaust air in extraction hoods. Smaller cyclones are used to separate airborne particles for analysis. Some are small enough to be worn clipped to clothing, and are used to separate respirable particles for later analysis.
Similar separators are used in the oil refining industry (e.g. for Fluid catalytic cracking) to achieve fast separation of the catalyst particles from the reacting gases and vapors.
Analogous devices for separating particles or solids from liquids are called hydrocyclones or hydroclones. These may be used to separate solid waste from water in wastewater and sewage treatment.
Types
The most common types of centrifugal, or inertial, collectors in use today are:
Single-cyclone separators
Single-cyclone separators create a dual vortex to separate coarse from fine dust. The main vortex spirals downward and carries most of the coarser dust particles. The inner vortex, created near the bottom of the cyclone, spirals upward and carries finer dust particles.
Multiple-cyclone separators
Multiple-cyclone separators consist of a number of small-diameter cyclones, operating in parallel and having a common gas inlet and outlet, as shown in the figure, and operate on the same principle as single cyclone separators—creating an outer downward vortex and an ascending inner vortex.
Multiple-cyclone separators remove more dust than single cyclone separators because the individual cyclones have a greater length and smaller diameter. The longer length provides longer residence time while the smaller diameter creates greater centrifugal force. These two factors result in better separation of dust particulates. The pressure drop of multiple-cyclone separators collectors is higher than that of single-cyclone separators, requiring more energy to clean the same amount of air. A single-chamber cyclone separator of the same volume is more economical, but doesn't remove as much dust.
Secondary-air-flow separators
This type of cyclone uses a secondary air flow, injected into the cyclone to accomplish several things. The secondary air flow increases the speed of the cyclonic action making the separator more efficient; it intercepts the particulate before it reaches the interior walls of the unit; and it forces the separated particulate toward the collection area. The secondary air flow protects the separator from particulate abrasion and allows the separator to be installed horizontally because gravity is not depended upon to move the separated particulate downward.
Cyclone theory
As the cyclone is essentially a two phase particle-fluid system, fluid mechanics and particle transport equations can be used to describe the behaviour of a cyclone. The air in a cyclone is initially introduced tangentially into the cyclone with an inlet velocity . Assuming that the particle is spherical, a simple analysis to calculate critical separation particle sizes can be established.
If one considers an isolated particle circling in the upper cylindrical component of the cyclone at a rotational radius of from the cyclone's central axis, the particle is therefore subjected to drag, centrifugal, and buoyant forces. Given that the fluid velocity is moving in a spiral the gas velocity can be broken into two component velocities: a tangential component, , and an outward radial velocity component . Assuming Stokes' law, the drag force in the outward radial direction that is opposing the outward velocity on any particle in the inlet stream is:
Using as the particle's density, the centrifugal component in the outward radial direction is:
The buoyant force component is in the inward radial direction. It is in the opposite direction to the particle's centrifugal force because it is on a volume of fluid that is missing compared to the surrounding fluid. Using for the density of the fluid, the buoyant force is:
In this case, is equal to the volume of the particle (as opposed to the velocity). Determining the outward radial motion of each particle is found by setting Newton's second law of motion equal to the sum of these forces:
To simplify this, we can assume the particle under consideration has reached "terminal velocity", i.e., that its acceleration is zero. This occurs when the radial velocity has caused enough drag force to counter the centrifugal and buoyancy forces. This simplification changes our equation to:
Which expands to:
Solving for we have
.
Notice that if the density of the fluid is greater than the density of the particle, the motion is (-), toward the center of rotation and if the particle is denser than the fluid, the motion is (+), away from the center. In most cases, this solution is used as guidance in designing a separator, while actual performance is evaluated and modified empirically.
In non-equilibrium conditions when radial acceleration is not zero, the general equation from above must be solved. Rearranging terms we obtain
Since is distance per time, this is a 2nd order differential equation of the form .
Experimentally it is found that the velocity component of rotational flow is proportional to , therefore:
This means that the established feed velocity controls the vortex rate inside the cyclone, and the velocity at an arbitrary radius is therefore:
Subsequently, given a value for , possibly based upon the injection angle, and a cutoff radius, a characteristic particle filtering radius can be estimated, above which particles will be removed from the gas stream.
Alternative models
The above equations are limited in many regards. For example, the geometry of the separator is not considered, the particles are assumed to achieve a steady state and the effect of the vortex inversion at the base of the cyclone is also ignored, all behaviours which are unlikely to be achieved in a cyclone at real operating conditions.
More complete models exist, as many authors have studied the behaviour of cyclone separators. Simplified models allowing a quick calculation of the cyclone, with some limitations, have been developed for common applications in process industries. Numerical modelling using computational fluid dynamics has also been used extensively in the study of cyclonic behaviour. A major limitation of any fluid mechanics model for cyclone separators is the inability to predict the agglomeration of fine particles with larger particles, which has a great impact on cyclone collection efficiency.
| Physical sciences | Other separations | Chemistry |
767304 | https://en.wikipedia.org/wiki/Camellia%20sinensis | Camellia sinensis | Camellia sinensis is a species of evergreen shrub or small tree in the flowering plant family Theaceae. Its leaves, leaf buds, and stems can be used to produce tea. Common names include tea plant, tea shrub, and tea tree (unrelated to Melaleuca alternifolia, the source of tea tree oil, or the genus Leptospermum commonly called tea tree).
White tea, yellow tea, green tea, oolong, dark tea (which includes pu-erh tea) and black tea are all harvested from one of two major varieties grown today, C. sinensis var. sinensis and C. s. var. assamica, but are processed differently to attain varying levels of oxidation with black tea being the most oxidized and white being the least. Kukicha (twig tea) is also harvested from C. sinensis, but uses twigs and stems rather than leaves.
Description
Camellia sinensis is native to East Asia, the Indian Subcontinent, and Southeast Asia, but it is today cultivated all around the world in tropical and subtropical regions. It is an evergreen shrub or small tree that is usually trimmed to below when cultivated for its leaves. It has a strong taproot. The flowers are yellow-white, in diameter, with seven or eight petals.
The seeds of C. sinensis and C. oleifera can be pressed to yield tea oil, a sweetish seasoning and cooking oil that should not be confused with tea tree oil, an essential oil that is used for medical and cosmetic purposes, and originates from the leaves of a different plant.
The leaves are long and broad. Fresh leaves contain about 4% caffeine, as well as related compounds including theobromine. The young, light-green leaves are preferably harvested for tea production when they have short, white hairs on the underside. Older leaves are deeper green. Different leaf ages produce differing tea qualities, since their chemical compositions are different. Usually, the tip (bud) and the first two to three leaves are harvested for processing. This hand picking is repeated every one to two weeks.
In 2017, Chinese scientists sequenced the genome of C. s. var. assamica. It contains about three billion base pairs, which was larger than most plants previously sequenced.
Taxonomy
The generic name Camellia is taken from the Latinized name of Rev. Georg Kamel, SJ (1661–1706), a Moravian-born Jesuit lay brother, pharmacist, and missionary to the Philippines.
Carl Linnaeus chose his name in 1753 for the genus to honor Kamel's contributions to botany (although Kamel did not discover or name this plant, or any Camellia, and Linnaeus did not consider this plant a Camellia but a Thea).
Robert Sweet shifted all formerly Thea species to the genus Camellia in 1818. The name sinensis means "from China" in Latin.
Four varieties of C. sinensis are recognized. Of these, C. sinensis var. sinensis and C. s. var. assamica (JW Masters) Kitamura are most commonly used for tea, and C. s. var. pubilimba Hung T. Chang and C. s. var. dehungensis (Hung T. Chang & BH Chen) TL Ming are sometimes used locally. The Cambodia type tea (C. assamica subsp. lasiocaly) was originally considered a type of assam tea. However, later genetic work showed that it is a hybrid between Chinese small leaf tea and assam type tea.
Tea plants are native to East Asia, and probably originated in the borderlands of north Burma and southwestern China.
Chinese (small leaf) tea [C. sinensis var. sinensis]
Chinese Western Yunnan Assam (large leaf) tea [C. sinensis var. assamica]
Indian Assam (large leaf) tea [C. sinensis var. assamica]
Chinese Southern Yunnan Assam (large leaf) tea [C. sinensis var. assamica]
Chinese (small leaf) tea may have originated in southern China possibly with hybridization of unknown wild tea relatives. However, since no wild populations of this tea are known, the precise location of its origin is speculative.
Given their genetic differences forming distinct clades, Chinese Assam type tea (C. s. var. assamica) may have two different parentages – one being found in southern Yunnan (Xishuangbanna, Pu'er City) and the other in western Yunnan (Lincang, Baoshan). Many types of Southern Yunnan Assam tea have been hybridized with the closely related species Camellia taliensis. Unlike Southern Yunnan Assam tea, Western Yunnan Assam tea shares many genetic similarities with Indian Assam type tea (also C. s. var. assamica). Thus, Western Yunnan Assam tea and Indian Assam tea both may have originated from the same parent plant in the area where southwestern China, Indo-Burma, and Tibet meet. However, as the Indian Assam tea shares no haplotypes with Western Yunnan Assam tea, Indian Assam tea is likely to have originated from an independent domestication. Some Indian Assam tea appears to have hybridized with the species Camellia pubicosta.
Assuming a generation of 12 years, Chinese small leaf tea is estimated to have diverged from Assam tea around 22,000 years ago; this divergence would correspond to the last glacial maximum, while Chinese Assam tea and Indian Assam tea diverged 2,800 years ago.
Chinese small leaf type tea was introduced into India in 1836 by the British and some Indian Assam type tea (e.g. Darjeeling tea) appear to be genetic hybrids of Chinese small leaf type tea, native Indian Assam, and possibly also closely related wild tea species.
Cultivars
Hundreds, if not thousands of cultivars of C. sinensis are known. Some Japanese cultivars include:
Benifuuki
Fushun
Kanayamidori
Meiryoku
Saemidori
Okumidori
Yabukita
Cultivation
Camellia sinensis is mainly cultivated in tropical and subtropical climates, in areas with at least 127 cm (50 in) of rainfall a year. Tea plants prefer a rich and moist growing location in full to part sun, and can be grown in hardiness zones 7–9. However, the clonal one is commercially cultivated from the equator to as far north as Cornwall and Scotland on the UK mainland. Many high quality teas are grown at high elevations, up to , as the plants grow more slowly and acquire more flavor.
Tea plants will grow into a tree if left undisturbed, but cultivated plants are pruned to waist height for ease of plucking. Two principal varieties are used, the small-leaved Chinese variety plant (C. s. sinensis) and the large-leaved Assamese plant (C. s. assamica), used mainly for black tea.
Chinese teas
The Chinese plant is a small-leafed bush with multiple stems that reaches a height of some . It is native to southeast China. The first tea plant variety to be discovered, recorded, and used to produce tea dates back 3,000 years ago; it yields some of the most popular teas.
C. s. var. waldenae was considered a different species, C. waldenae by SY Hu, but it was later identified as a variety of C. sinensis. This variety is commonly called Waldenae Camellia. It is seen on Sunset Peak and Tai Mo Shan in Hong Kong. It is also distributed in the Guangxi province.
Indian teas
Three main kinds of tea are produced in India:
Assam, from the var. assamica plant, comes from the near sea-level heavily forested northeastern section of India, the state of Assam. Tea from here is rich and full-bodied. The first tea estate in India was established in Assam in 1837. Teas are manufactured in either the orthodox process or the "crush, tear, curl" (CTC) process.
Darjeeling, from the var. sinensis plant, is from the cool and wet Darjeeling highland region, tucked in the foothills of the Himalayas. Tea plantations could be at altitudes as high as . The tea is delicately flavored, and considered to be one of the finest teas in the world. The Darjeeling plantations have three distinct harvests, termed 'flushes', and the tea produced from each flush has a unique flavor. First (spring) flush teas are light and aromatic, while the second (summer) flush produces tea with a bit more bite. The third, or autumn flush gives a tea that is lesser in quality.
Nilgiri is from a southern region of India almost as high as Darjeeling. Grown at elevations between , Nilgiri teas are subtle and rather gentle, and are frequently blended with other, more robust teas.
Pests and diseases
Tea leaves are eaten by some herbivores, such as the caterpillars of the willow beauty (Peribatodes rhomboidaria), a geometer moth.
Health effects
Although health benefits have been assumed throughout the history of using tea as a common beverage, no high-quality evidence shows that tea confers significant benefits. In clinical research over the early 21st century, tea has been studied extensively for its potential to lower the risk of human diseases, but none of this research is conclusive as of 2017.
Biosynthesis of caffeine
Caffeine, a molecule produced in C. sinensis, functions as a secondary metabolite and acts as a natural pesticide: it can paralyze and kill herbivorous insects feeding on the plant. Caffeine is a purine alkaloid and its biosynthesis occurs in young tea leaves and is regulated by several enzymes. The biosynthetic pathway in C. sinensis is similar to other caffeine-producing plants such as coffee or guayusa. Analysis of the pathway was carried out by harvesting young leaves and using reverse transcription PCR to analyze the genes encoding the major enzymes involved in synthesizing caffeine. The gene TCS1 encodes caffeine synthase. Younger leaves feature high concentrations of TCS1 transcripts, allowing more caffeine to be synthesized during this time. Dephosphorylation of xanthosine-5'-monophosphate into xanthosine is the committed step for the xanthosines entering the beginning of the most common pathway. A sequence of reactions turns xanthosine (9β--ribofuranosylxanthine) into 7-methylxanthosine, then 7-methylxanthine, then theobromine (3,7-dimethylxanthine), and finally into caffeine (1,3,7-trimethylxanthine).
| Biology and health sciences | Ericales | Plants |
767894 | https://en.wikipedia.org/wiki/Pocket%20watch | Pocket watch | A pocket watch is a watch that is made to be carried in a pocket, as opposed to a wristwatch, which is strapped to the wrist.
They were the most common type of watch from their development in the 16th century until wristwatches became popular after World War I during which a transitional design, trench watches, were used by the military. Pocket watches generally have an attached chain to allow them to be secured to a waistcoat, lapel, or belt loop, and to prevent them from being dropped. Watches were also mounted on a short leather strap or fob, when a long chain would have been cumbersome or likely to catch on things. This fob could also provide a protective flap over their face and crystal. Women's watches were normally of this form, with a watch fob that was more decorative than protective. Chains were frequently decorated with a silver or enamel pendant, often carrying the arms of some club or society, which by association also became known as a fob. Ostensibly practical gadgets such as a watch winding key, vesta case, seal, and/or a cigar cutter also appeared on watch chains, although usually in an overly decorated style. Also common are fasteners designed to be put through a buttonhole and worn in a jacket or waistcoat, this sort being frequently associated with and named after train conductors.
An early reference to the pocket watch is in a letter in November 1462 from the Italian clockmaker Bartholomew Manfredi to the Marchese di Mantova Federico Gonzaga, where he offers him a "pocket clock" better than that belonging to the Duke of Modena. By the end of the 15th century, spring-driven clocks appeared in Italy, and in Germany. Peter Henlein, a master locksmith of Nuremberg, was regularly manufacturing pocket watches by 1526. Thereafter, pocket watch manufacture spread throughout the rest of Europe as the 16th century progressed. Early watches only had an hour hand, the minute hand appearing in the late 17th century.
History
The first timepieces to be worn, made in 16th-century Europe, were transitional in size between clocks and watches. These 'clock-watches' were fastened to clothing or worn on a chain around the neck. They were heavy drum shaped brass cylinders several inches in diameter, engraved and ornamented. They had only an hour hand. The face was not covered with glass, but usually had a hinged brass cover, often decoratively pierced with grillwork so the time could be read without opening. The movement was made of iron or steel and held together with tapered pins and wedges, until screws began to be used after 1550. Many of the movements included striking or alarm mechanisms. The shape later evolved into a rounded form; these were later called Nuremberg eggs. Still later in the century there was a trend for unusually shaped watches, and clock-watches shaped like books, animals, fruit, stars, flowers, insects, crosses, and even skulls (Death's head watches) were made.
Styles changed in the 17th century and men began to wear watches in pockets instead of as pendants (the woman's watch remained a pendant into the 20th century). This is said to have occurred in 1675 when Charles II of England introduced waistcoats. To fit in pockets, their shape evolved into the typical pocket watch shape, rounded and flattened with no sharp edges. Glass was used to cover the face beginning around 1610. Watch fobs began to be used, the name originating from the German word fuppe, a small pocket. The watch was wound and also set by opening the back and fitting a key to a square arbor, and turning it.
Until the second half of the 18th century, watches were luxury items; as an indication of how highly they were valued, English newspapers of the 18th century often include advertisements offering rewards of between one and five guineas merely for information that might lead to the recovery of stolen watches. By the end of the 18th century, however, watches (while still largely hand-made) were becoming more common; special cheap watches were made for sale to sailors, with crude but colorful paintings of maritime scenes on the dials.
Up to the 1720s, almost all watch movements were based on the verge escapement, which had been developed for large public clocks in the 14th century. This type of escapement involved a high degree of friction and did not include any kind of jewelling to protect the contacting surfaces from wear. As a result, a verge watch could rarely achieve any high standard of accuracy. (Surviving examples mostly run very fast, often gaining an hour a day or more.) The first widely used improvement was the cylinder escapement, developed by the Abbé de Hautefeuille early in the 18th century and applied by the English maker George Graham. Then, towards the end of the 18th century, the lever escapement (invented by Thomas Mudge in 1755) was put into limited production by a handful of makers including Josiah Emery (a Swiss based in London) and Abraham-Louis Breguet. With this, a domestic watch could keep time to within a minute a day. Lever watches became common after about 1820, and this type is still used in most mechanical watches.
In 1857 the American Watch Company in Waltham, Massachusetts, introduced the Waltham Model 57, the first to use interchangeable parts. This cut the cost of manufacture and repair. Most Model 57 pocket watches were in a coin silver ("one nine fine"), a 90% pure silver alloy commonly used in dollar coinage, slightly less pure than the British (92.5%) sterling silver, both of which avoided the higher purity of other types of silver to make circulating coins and other utilitarian silver objects last longer with heavy use.
Watch manufacture was becoming streamlined; the Japy family of Schaffhausen, Switzerland, led the way in this, and soon afterwards the newborn American watch industry developed much new machinery, so that by 1865 the American Watch Company (afterwards known as Waltham) could turn out more than 50,000 reliable watches each year. This development drove the Swiss out of their dominating position at the cheaper end of the market, compelling them to raise the quality of their products and establish themselves as the leaders in precision and accuracy instead.
Use in railroading in the United States
The rise of railroading during the last half of the 19th century led to the widespread use of pocket watches.
A famous train wreck on the Lake Shore and Michigan Southern Railway in Kipton, Ohio, on April 19, 1891, occurred because one of the engineers' watches had stopped for four minutes. The railroad officials commissioned Webb C. Ball as their Chief Time Inspector, in order to establish precision standards and a reliable timepiece inspection system for Railroad chronometers. This led to the adoption in 1893 of stringent standards for pocket watches used in railroading. These railroad-grade pocket watches, as they became colloquially known, had to meet the General Railroad Timepiece Standards adopted in 1893 by almost all railroads. These standards read, in part:
Types of pocket watches
There are two main styles of pocket watch, the hunter-case pocket watch and the open-face pocket watch.
Open-face watches
An open-faced, or Lépine, watch, is one in which the case lacks a metal cover to protect the crystal. It is typical for an open-faced watch to have the pendant located at 12:00 and the sub-second dial located at 6:00. Occasionally, a watch movement intended for a hunting case (with the winding stem at 3:00 and sub second dial at 6:00) will have an open-faced case. Such watch is known as a "sidewinder." Alternatively, such a watch movement may be fitted with a so-called conversion dial, which relocates the winding stem to 12:00 and the sub-second dial to 3:00. After 1908, watches approved for railroad service were required to be cased in open-faced cases with the winding stem at 12:00.
Hunter-case watches
A hunter-case pocket watch is a case with a spring-hinged circular metal lid or cover, that closes over the watch-dial and crystal, protecting them from dust, scratches and other damage or debris. The name originated from England where "fox hunting men found it convenient to be able to open their watch and read the time with one hand, while holding the reins of their 'hunter' (horse) in the other hand".
It is also known as a "savonnette", after the French word for soap (savon) due to its resemblance to a round soap bar.
The majority of antique and vintage hunter-case watches have the lid-hinges at the 9 o'clock position and the stem, crown and bow of the watch at the 3 o'clock position. Modern hunter-case pocket watches usually have the hinges for the lid at the 6 o'clock position and the stem, crown and bow at the 12 o'clock position, as with open-face watches. In both styles of watch-cases, the sub-seconds dial was always at the 6 o'clock position. A hunter-case pocket watch with a spring-ring chain is pictured at the top of this page.
An intermediate type, known as the demi-hunter (or half-hunter), is a case style in which the outer lid has a glass panel or hole in the centre giving a view of the hands. The hours are marked, often in blue enamel, on the outer lid itself; thus with this type of case one can tell the time without opening the lid.
Types of watch movements
Key-wind, key-set movements
The first pocket watches, since their creation in the 16th century, up until the third quarter of the 19th century, had key-wind and key-set movements. A watch key was necessary to wind the watch and to set the time. This was usually done by opening the caseback and putting the key over the winding-arbor (which was set over the watch's winding-wheel, to wind the mainspring) or by putting the key onto the setting-arbor, which was connected with the minute-wheel and turned the hands. Some watches of this period had the setting-arbor at the front of the watch, so that removing the crystal and bezel was necessary to set the time. Watch keys are the origin of the class key, common paraphernalia for American high-school and university graduation.
Many keywind watch movements make use of a fusee, to improve isochronism. The fusee is a specially cut conical pulley attached by a fine chain to the mainspring barrel. When the spring is fully wound (and its torque the highest), the full length of the chain is wrapped around the fusee and the force of the mainspring is exerted on the smallest diameter portion of the fusee cone. As the spring unwinds and its torque decreases, the chain winds back onto the mainspring barrel and pulls on an increasingly larger diameter portion of the fusee. This provides a more uniform amount of torque on the watch train, and thus results in more consistent balance amplitude and better isochronism. A fusee is a practical necessity in watches using a verge escapement, and can also provide considerable benefit with a lever escapement and other high precision types of escapements (Hamiltons WWII era Model 21 chronometer used a fusee in combination with a detent escapement).
Keywind watches are also commonly seen with conventional going barrels and other types of mainspring barrels, particularly in American watchmaking.
Stem-wind, stem-set movements
Invented by Adrien Philippe in 1842 and commercialized by Patek Philippe & Co. in the 1850s, the stem-wind, stem-set movement did away with the watch key which was a necessity for the operation of any pocket watch up to that point. The first stem-wind and stem-set pocket watches were sold during the Great Exhibition in London in 1851 and the first owners of these new kinds of watches were Queen Victoria and Prince Albert. Stem-wind, stem-set movements are the most common type of watch-movement found in both vintage and modern pocket watches.
The mainstream transition to the use of stem-wind, stem-set watches occurred at around the same time as the end of the manufacture and use of the fusee watch. Fusee chain-driven timing was replaced with a mainspring of better quality spring steel (commonly known as the "going barrel") allowing for a more even release of power to the escape mechanism. The balance wheel and balance spring provide a separate function: to regulate the timing (or escape) of the movement.
Stem-wind, lever-set movements
Mandatory for all railroad watches after roughly 1908, this kind of pocket watch was set by opening the crystal and bezel and pulling out the setting-lever (most hunter-cases have levers accessible without removing the crystal or bezel), which was generally found at either the 10 or 2 o'clock positions on open-faced watches, and at 5:00 on hunting cased watches. Once the lever was pulled out, the crown could be turned to set the time. The lever was then pushed back in and the crystal and bezel were closed over the dial again. This method of time setting on pocket watches was preferred by American and Canadian railroads, as lever setting watches make accidental time changes impossible. After 1908, lever setting was generally required for new watches entering service on American railroads.
Stem-wind, pin-set movements
Much like the lever-set movements, these pocket watches had a small pin or knob next to the watch-stem that had to be depressed before turning the crown to set the time and releasing the pin when the correct time had been set. This style of watch is occasionally referred to as "nail set", as the set button must be pressed using a fingernail.
Adjusted movements
Pocket watch movements are occasionally engraved with the word "Adjusted", or "Adjusted to n positions". This means that the watch has been tuned to keep time under various positions and conditions. There are eight possible adjustments:
Dial up
Dial down
Pendant up
Pendant down
Pendant left
Pendant right
Temperature (from )
Isochronism (the ability of the watch to keep time, regardless of the mainspring's level of tension)
Positional adjustments are attained by careful poising (ensuring even weight distribution) of the balance-hairspring system as well as careful control of the shape and polish on the balance pivots. All of this achieves an equalization of the effect of gravity on the watch in various positions. Positional adjustments are achieved through careful adjustment of each of these factors, provided by repeated trials on a timing machine. Thus, adjusting a watch to position requires many hours of labor, increasing the cost of the watch. Medium grade watches were commonly adjusted to 3 positions (dial up, dial down, pendant up) while high grade watches were commonly adjusted to 5 positions (dial up, dial down, stem up, stem left, stem right) or even all 6 positions. Railroad watches were required, after 1908, to be adjusted to 5 positions. 3 positions were the general requirement before that time.
Early watches used a solid steel balance. As temperature increased, the solid balance expanded in size, changing the moment of inertia and changing the timing of the watch. In addition, the hairspring would lengthen, decreasing its spring constant. This problem was initially solved through the use of the compensation balance. The compensation balance consisted of a ring of steel sandwiched to a ring of brass. These rings were then split in two places. The balance would, at least theoretically, actually decrease in size with heating to compensate for the lengthening of the hairspring. Through careful adjustment of the placement of the balance screws (brass or gold screws placed in the rim of the balance), a watch could be adjusted to keep time the same at both hot () and cold () temperatures. Unfortunately, a watch so adjusted would run slow at temperatures between these two. The problem was completely solved through the use of special alloys for the balance and hairspring which were essentially immune to thermal expansion. Such an alloy is used in Hamilton's 992E and 992B.
Isochronism was occasionally improved through the use of a stopworks, a system designed to only allow the mainspring to operate within its center (most consistent) range. The most common method of achieving isochronism is through the use of the Breguet overcoil, which places part of the outermost turn of the hairspring in a different plane from the rest of the spring. This allows the hairspring to "breathe" more evenly and symmetrically. Two types of overcoils are found - the gradual overcoil and the Z-Bend. The gradual overcoil is obtained by imposing two gradual twists to the hairspring, forming the rise to the second plane over half the circumference; and the Z-bend does this by imposing two kinks of complementary 45 degree angles, accomplishing a rise to the second plane in about three spring section heights. The second method is done for esthetic reasons and is much more difficult to perform. Due to the difficulty with forming an overcoil, modern watches often use a slightly less effective "dogleg", which uses a series of sharp bends (in plane) to place part of the outermost coil out of the way of the rest of the spring.
Decline in popularity
Pocket watches are uncommon in the present day, having first been superseded by wristwatches and later by smartphones. Until early in the 20th century, though, the pocket watch remained predominant for men, with the wristwatch considered feminine and unmanly. In men's fashions, pocket watches began to be superseded by wristwatches around the time of World War I, when officers in the field began to appreciate that a watch worn on the wrist was more easily accessed than one kept in a pocket. A watch of transitional design, combining features of pocket watches and modern wristwatches, was called a "trench watch" or "wristlet". The more accurate pocket watches continued to be widely used in railroading even as their popularity declined elsewhere. Quartz pocket watches are available in the present day, retaining the form and function of the original pocket watches while using a quartz crystal as opposed to the traditional fully-mechanical movement.
For a few years in the late 1970s and early 1980s three-piece suits for men returned to fashion, and this led to small resurgence in pocket watches, as some men actually used the vest pocket for its original purpose. Since then, some watch companies continue to make pocket watches. As vests have long since fallen out of fashion (in the US) as part of formal business wear, the only available location for carrying a watch is in a trouser pocket. The more recent advent of mobile phones and other gadgets that are worn on the waist has diminished the appeal of carrying an additional item in the same location, especially as such pocketable gadgets usually have timekeeping functionality themselves.
In some countries, namely the US, a gift of a gold-cased pocket watch is traditionally awarded to an employee upon their retirement.
The pocket watch has regained popularity with the steampunk subcultural movement embracing the arts and fashions of the Victorian era, during which pocket watches were nearly ubiquitous.
In animated films and video games, especially within the fantasy genre, devices resembling pocket watches commonly represent objects with the ability to alter time, such as by time travel.
Most complicated pocket watches
The Vacheron Constantin Reference 57260 (2015) – 57 complications
Patek Philippe Calibre 89 (1989) – 33 complications
Patek Philippe Henry Graves Supercomplication (1933) – 24 complications
| Technology | Clocks | null |
8119544 | https://en.wikipedia.org/wiki/Monosodium%20citrate | Monosodium citrate | Monosodium citrate, more correctly, sodium dihydrogen citrate (Latin: ), is an acid salt of citric acid. Disodium citrate and trisodium citrate are also known. It can be prepared by partial neutralisation of citric acid with an aqueous solution of sodium bicarbonate or carbonate. It has a slightly acidic taste.
NaHCO3 + C6H8O7 → NaC6H7O7 + CO2 + H2O
Na2CO3 + 2C6H8O7 → 2NaC6H7O7 + CO2 + H2O
It is highly soluble in water and practically insoluble in ethanol. Monosodium citrate is used as an anticoagulant in donated blood. It is used as an alkalinizing agent to prevent kidney stone disease. The crystals form as nearly perfect cubes.
| Physical sciences | Citrates | Chemistry |
8119545 | https://en.wikipedia.org/wiki/Disodium%20citrate | Disodium citrate | Disodium citrate, also known as disodium hydrogen citrate, (Neo-Alkacitron) and sesquihydrate, is an acid salt of citric acid with the chemical formula . It is used as an antioxidant in food and to improve the effects of other antioxidants. It is also used as an acidity regulator and sequestrant. Typical products include gelatin, jam, sweets, ice cream, carbonated beverages, milk powder, wine, and processed cheeses.
Uses
Food
It is used as an antioxidant in food and to improve the effects of other antioxidants. It is also used as an acidity regulator and sequestrant. Typical products include gelatin, jam, sweets, ice cream, carbonated beverages, milk powder, wine, and processed cheeses. Disodium citrate can also be used as a thickening agent or stabilizer.
Manufacturing
Disodium citrate can also be used as an ingredient in household products that remove stains.
Health
Disodium citrate may be used in patients to alleviate discomfort from urinary-tract infections.
| Physical sciences | Citrates | Chemistry |
116035 | https://en.wikipedia.org/wiki/Tetraethyllead | Tetraethyllead | Tetraethyllead (commonly styled tetraethyl lead), abbreviated TEL, is an organolead compound with the formula Pb(C2H5)4. It was widely used as a fuel additive for much of the 20th century, first being mixed with gasoline beginning in the 1920s. This "leaded gasoline" had an increased octane rating that allowed engine compression to be raised substantially and in turn increased vehicle performance and fuel economy. TEL was first synthesised by German chemist Carl Jacob Löwig in 1853. American chemical engineer Thomas Midgley Jr., who was working for the U.S. corporation General Motors, was the first to discover its effectiveness as an antiknock agent in 1921, after spending several years attempting to find an additive that was both highly effective and inexpensive.
In the mid-20th century, scientists discovered TEL caused lead poisoning and was toxic to the human brain, especially in children. The United States and many other countries began phasing out the use of TEL in automotive fuel in the 1970s. By the early 2000s, most countries had completely banned the use of TEL in gasoline. Since 2011, leaded gasoline has been banned in every country. A 2011 study backed by the United Nations estimated that the removal of TEL had resulted in $2.4 trillion in annual benefits, and 1.2 million fewer premature deaths.
Despite being banned from use in automotive fuel, TEL is sometimes still used in certain grades of aviation fuel. Innospec has claimed to be the last firm legally making TEL but, , TEL was being produced illegally by several companies in China. In July 2021, the sale of leaded gasoline for cars was completely phased out worldwide, prompting the United Nations Environment Programme (UNEP) to declare an "official end" of its use in cars on August 30, 2021.
Synthesis and properties
TEL is produced by reacting chloroethane with a sodium–lead alloy.
The product is recovered by steam distillation, leaving a sludge of lead and sodium chloride. TEL is a viscous colorless liquid with a sweet odor. Because TEL is charge neutral and contains an exterior of alkyl groups, it is highly lipophilic and soluble in petrol (gasoline). This property, which allows it to dissolve so evenly and effectively in motor fuel, also allows it to dissolve oils and fats well, and therefore, diffuse through the blood–brain barrier and accumulate within the limbic forebrain, frontal cortex, and hippocampus.
Despite decades of research, no reactions were found to improve upon this process, which is rather difficult, involves metallic sodium, and converts only 25% of the lead to TEL. A related compound, tetramethyllead, was commercially produced by a different electrolytic reaction. A process with lithium was developed but never put into practice.
Reactions
A noteworthy feature of TEL is the weakness of its four C–Pb bonds. At the temperatures found in internal combustion engines, TEL decomposes completely into lead as well as combustible, short-lived ethyl radicals. Lead and lead oxide scavenge radical intermediates in combustion reactions. Engine knock is caused by a cool flame, an oscillating low-temperature combustion reaction that occurs before the proper, hot ignition. Lead quenches the pyrolysed radicals and thus kills the radical chain reaction that would sustain a cool flame, preventing it from disturbing the smooth ignition of the hot flame front. Lead itself is the reactive antiknock agent, and the ethyl groups serve as a gasoline-soluble carrier.
When TEL burns, it produces not only carbon dioxide and water, but also lead and lead(II) oxide:
Pb and PbO would quickly over-accumulate and foul an engine. For this reason, 1,2-dichloroethane and 1,2-dibromoethane were also added to gasoline as lead scavengers—these agents form volatile lead(II) chloride and lead(II) bromide, respectively, which flush the lead from the engine and into the air:
In motor fuel
TEL was extensively used as a gasoline additive beginning in the 1920s, wherein it served as an effective antiknock agent and reduced exhaust valve and valve seat wear. Concerns were raised in reputable journals of likely health outcomes of fine particles of lead in the atmosphere as early as 1924.
Valve wear preventive
Tetraethyllead helps cool intake valves and is an excellent buffer against microwelds forming between exhaust valves and their seats. Once these valves reopen, the microwelds pull apart and abrade the valves and seats, leading to valve recession. When TEL began to be phased out, the automotive industry began specifying hardened valve seats and upgraded materials which allow for high wear resistance without requiring lead.
Antiknock agent
A gasoline-fuelled reciprocating engine requires fuel of sufficient octane rating to prevent uncontrolled combustion (preignition and detonation). Antiknock agents allow the use of higher compression ratios for greater efficiency and peak power. Adding varying amounts of additives to gasoline allowed easy, inexpensive control of octane ratings. TEL offered the business advantage of being commercially profitable because its use for this purpose could be patented. Aviation fuels with TEL used in WWII reached octane ratings of 150 to enable turbocharged and supercharged engines such as the Rolls-Royce Merlin and Griffon to reach high horsepower ratings at altitude. In military aviation, TEL manipulation allowed a range of different fuels to be tailored for particular flight conditions.
In 1935 a licence to produce TEL was given to IG Farben, enabling the newly formed German to use high-octane gasoline. A company, Ethyl GmbH, was formed that produced TEL at two sites in Germany with a government contract from 10 June 1936.
In 1938 the United Kingdom Air Ministry contracted with ICI for the construction and operation of a TEL plant. A site was chosen at Holford Moss, near Plumley in Cheshire. Construction started in April 1939 and TEL was being produced by September 1940.
"Ethyl Fluid"
For mixing with raw gasoline, TEL was most commonly supplied in the form of "Ethyl Fluid", which consisted of TEL blended with 1,2-dichloroethane and 1,2-dibromoethane. Ethyl Fluid also contained a reddish dye to distinguish treated from untreated gasoline and discourage the use of leaded gasoline for other purposes such as cleaning.
In the 1920s before safety procedures were strengthened, 17 workers for the Ethyl Corporation, DuPont, and Standard Oil died from the effects of exposure to lead.
Ethyl Fluid's formulation consisted of:
61.45% tetraethyllead
18.80% 1,2-dichloroethane
17.85% 1,2-dibromoethane
1.90% inerts and dyes
Dichloroethane and dibromoethane act in a synergistic manner, where equal or approximately equal quantities of both provide the best scavenging ability.
Phaseout and ban
In most industrialized countries, a phaseout of TEL from road vehicle fuels was completed by the early 2000s because of concerns over air and soil lead levels and the accumulative neurotoxicity of lead. In the European Union, tetraethyllead has been classified as a Substance of Very High Concern and placed on the Candidate List for Authorisation under Registration, Evaluation, Authorisation and Restriction of Chemicals (REACH). Potential use of TEL would need to be authorised through the REACH authorisation procedure. While not a complete ban, it introduces significant obligations such as a mandatory analysis of alternatives and socioeconomic analysis.
The use of catalytic converters, mandated in the United States for 1975 and later model-year cars to meet tighter emissions regulations, started a gradual phase-out of leaded gasoline in the U.S. The need for TEL was lessened by several advances in automotive engineering and petroleum chemistry. Safer methods for making higher-octane blending stocks such as reformate and iso-octane reduced the need to rely on TEL, as did other antiknock additives of varying toxicity including metallic compounds such as methylcyclopentadienyl manganese tricarbonyl (MMT) as well as oxygenates including methyl tert-butyl ether (MTBE), tert-amyl methyl ether (TAME), and ethyl tert-butyl ether (ETBE).
The first country to completely ban leaded gasoline was Japan in 1986.
Since January 1993 all gasoline powered cars sold in the European Union and the United Kingdom have been required to use unleaded fuel. This was to comply with the Euro 1 emission standards which mandated that all new cars to be fitted with a catalytic converter. Unleaded fuel was first introduced in the United Kingdom in June 1986.
Leaded gasoline was removed from the forecourts in the United Kingdom on January 1, 2000, and a Lead Replacement Petrol was introduced although this was largely withdrawn by 2003 due to dwindling sales. An exemption to the ban exists for owners of classic cars.
Vehicles designed and built to run on leaded fuel often require modification to run on unleaded gasoline. These modifications fall into two categories: those required for physical compatibility with unleaded fuel, and those performed to compensate for the relatively low octane of early unleaded fuels. Physical compatibility requires the installation of hardened exhaust valves and seats. Compatibility with reduced octane was addressed by reducing compression, generally by installing thicker cylinder head gaskets and/or rebuilding the engine with compression-reducing pistons, although modern high-octane unleaded gasoline has eliminated the need to decrease compression ratios.
Leaded gasoline remained legal as of late 2014 in parts of Algeria, Iraq, Yemen, Myanmar, North Korea, and Afghanistan. North Korea and Myanmar purchased their TEL from China, while Algeria, Iraq, and Yemen purchased it from the specialty chemical company Innospec, the world's sole remaining legal manufacturer of TEL. In 2011 several Innospec executives were charged and imprisoned for bribing various government state-owned oil companies to approve the sale of their TEL products.
the UNEP-sponsored phase-out was nearly complete: only Algeria, Iraq, and Yemen continued widespread use of leaded gasoline, although not exclusively. In July 2021, Algeria had halted its sale.
Leaded-fuel bans
Leaded-fuel bans for road vehicles came into effect as follows:
Europe
Armenia: 2001
Austria: 1989
Belarus: 1998
Bulgaria: 2002
Bosnia and Herzegovina: 2009
Croatia: 2006
Cyprus: 2004
Czech Republic: 2001
Denmark: 1994
European Union: 1 January 2000
Finland: 1994
France: 2000
Germany: 1996
Gibraltar: 2001
Greece: 2002
Hungary: 1999
Ireland: 1 January 2000
Italy: 1 January 2002
Malta: 2003
Monaco: 2000
Netherlands: 1998
Norway: 1997
Poland: December 2000
Slovenia: 2001
Spain: 1 August 2001
Portugal: 1999
Romania: 2005
Russia: 2003
Serbia: 2010
Sweden: 1995
Switzerland: 2000
Ukraine: 2003
United Kingdom: 1 January 2000
North America
Anguilla: 1998
Antigua and Barbuda: 1991
Aruba: 1997
Bahamas: 1996
Belize: 1997
Bermuda: 1990
Cayman Islands: 1999
Canada: December 1990
Costa Rica: 1996
Dominican Republic: 1999
El Salvador: 1992
Guatemala: 1991
Haiti: 1998
Honduras: 1996
Jamaica: 2000
Mexico: 1998
Nicaragua: 1995
Panama: 2002
Trinidad and Tobago: 2000
United States (including Puerto Rico): 1 January 1996
California: 1992
South America
Argentina: 1998
Bolivia: 1995
Brazil: 1989 or 1991
Chile: 2001 or 2005
Colombia: 1991
Guyana: 2000
Peru: 2004
Suriname: 2001
Uruguay: 2004
Venezuela: 2005
Asia
Afghanistan: 2016
Bangladesh: 1999
China: 2000
Hong Kong: 1999
India: March 2000
Saudi Arabia: 2001
Indonesia: 2006
Iran: 2003
Iraq: 2018
Japan: 1986
Malaysia: 2000
Myanmar: 2016
Nepal: 2000
North Korea: 2016
Pakistan: 2001
Philippines: 2000
Singapore: 1998
South Korea: 1993
Sri Lanka: 1999
Taiwan: 2000
Thailand: 1996
Turkey: 2006
United Arab Emirates: 2003
Vietnam: 2001
Yemen: 2018
Oceania
Australia: 2002
New Zealand: 1996
Guam: 1 January 1996 (USA)
Samoa: 2001
Africa
Egypt: 1999
South Africa: 2006
Leaded petrol was supposed to be completely phased out continent-wide on 1 January 2006, following a ban initiated from the 2002 Earth Summit. However, in Algeria refineries needed to be altered; as a result, leaded fuel remained available in parts of Algeria, with phaseout scheduled for 2016. After the Algerian Government outlawed the sale of leaded petrol throughout all of Algeria, leaded petrol has now been effectively phased out.
In motor racing
Leaded fuel was commonly used in professional motor racing, until its phase out beginning in the 1990s. Since 1992, Formula One racing cars have been required to use fuel containing no more than 5 mg/L of lead.
NASCAR began experimentation in 1998 with an unleaded fuel, and in 2006 began switching the national series to unleaded fuel, completing the transition at the Fontana round in February 2007 when the premier class switched. This was influenced after blood tests of NASCAR teams revealed elevated blood lead levels.
Aviation gasoline
TEL remains an ingredient of 100 octane avgas for piston-engine aircraft. The current formulation of 100LL (low lead, blue) aviation gasoline contains of TEL, half the amount of the previous 100/130 (green) octane avgas (at 4.24 grams per gallon), and twice as much as the 1 gram per gallon permitted in regular automotive leaded gasoline prior to 1988 and substantially greater than the allowed 0.001 grams per gallon in automotive unleaded gasoline sold in the United States today. The United States Environmental Protection Agency, FAA, and others are working on economically feasible replacements for leaded avgas, which still releases 100 tons of lead every year.
Alternative antiknock agents
Antiknock agents are classed as high-percentage additives, such as alcohol, and low-percentage additives based on heavy elements. Since the main problem with TEL is its lead content, many alternative additives that contain less poisonous metals have been examined. A manganese-carrying additive, methylcyclopentadienyl manganese tricarbonyl (MMT or methylcymantrene), was used for a time as an antiknock agent, though its safety is controversial and it has been the subject of bans and lawsuits. Ferrocene, an organometallic compound of iron, is also used as an antiknock agent although with some significant drawbacks.
High-percentage additives are organic compounds that do not contain metals, but require much higher blending ratios, such as 20–30% for benzene and ethanol. It had been established by 1921 that ethanol was an effective antiknock agent, but TEL was introduced instead mainly for commercial reasons. Oxygenates such as TAME derived from natural gas, MTBE made from methanol, and ethanol-derived ETBE, have largely supplanted TEL. MTBE has environmental risks of its own and there are also bans on its use.
Improvements to gasoline itself decrease the need for antiknock additives. Synthetic iso-octane and alkylate are examples of such blending stocks. Benzene and other high-octane aromatics can be also blended to raise the octane number, but they are disfavored today because of toxicity and carcinogenicity.
Toxicity
6 mL of tetraethyllead is enough to induce severe lead poisoning. The hazards of TEL content are heightened due to the compound's volatility and high lipophilicity, enabling it to easily cross the blood–brain barrier.
Early symptoms of acute exposure to tetraethyllead can manifest as irritation of the eyes and skin, sneezing, fever, vomiting, and a metallic taste in the mouth. Later symptoms of acute TEL poisoning include pulmonary edema, anemia, ataxia, convulsions, severe weight loss, delirium, irritability, hallucinations, nightmares, fever, muscle and joint pain, swelling of the brain, coma, and damage to cardiovascular and renal organs.
Chronic exposure to TEL can cause long-term negative effects such as memory loss, delayed reflexes, neurological problems, insomnia, tremors, psychosis, loss of attention, and an overall decrease in IQ and cognitive function.
The carcinogenity of tetraethyllead is debatable. It is believed to harm the male reproductive system and cause birth defects.
Concerns over the toxicity of lead eventually led to the ban on TEL in automobile gasoline in many countries. Some neurologists have speculated that the lead phaseout may have caused average IQ levels to rise by several points in the US (by reducing cumulative brain damage throughout the population, especially in the young). For the entire US population, during and after the TEL phaseout, the mean blood lead level dropped from 16 μg/dL in 1976 to only 3 μg/dL in 1991. The U.S. Centers of Disease control previously labelled children with 10 μg/dL or more as having a "blood lead level of concern". In 2021, the level was lowered in accordance with the average lead level in the U.S. decreasing to 3.5 μg/dL or more as having a "blood lead level of concern".
History
In 1853, German chemist Karl Jacob Löwig (1803–1890) first prepared what he claimed was Pb2(C2H5)3 from ethyl iodide and an alloy of lead and sodium. In 1859, English chemist George Bowdler Buckton (1818–1905) reported what he claimed was Pb(C2H5)2 from zinc ethyl (Zn(C2H5)2) and lead(II) chloride. Later authors credit both methods of preparation with producing tetraethyl lead.
In fuel
TEL remained unimportant commercially until the 1920s. In 1921, at the direction of DuPont Corporation, which manufactured TEL, it was found to be an effective antiknock agent by Thomas Midgley, working under Charles Kettering at General Motors Corporation Research. General Motors patented the use of TEL as an antiknock agent and used the name "Ethyl" that had been proposed by Kettering in its marketing materials, thereby avoiding the negative connotation of the word "lead". Early research into "engine knocking" (also called "pinging" or "pinking") was also led by A.H. Gibson and Harry Ricardo in England and Thomas Boyd in the United States. The discovery that lead additives modified this behavior led to the widespread adoption of their use in the 1920s, and therefore more powerful, higher-compression engines. In 1924, Standard Oil of New Jersey (ESSO/EXXON) and General Motors created the Ethyl Gasoline Corporation to produce and market TEL. Deepwater, New Jersey, across the river from Wilmington, was the site for production of some of DuPont's most important chemicals, particularly TEL. After TEL production at the Bayway Refinery was shut down, Deepwater was the only plant in the Western hemisphere producing TEL up to 1948, when it accounted for the bulk of the Dupont/Deepwater's production.
Initial controversy
The toxicity of concentrated TEL was recognized early on, as lead had been recognized since the 19th century as a dangerous substance that could cause lead poisoning. In 1924, a public controversy arose over the "loony gas", after five workers died, and many others were severely injured, in Standard Oil refineries in New Jersey. There had also been a private controversy for two years prior to this controversy; several public health experts, including Alice Hamilton and Yandell Henderson, engaged Midgley and Kettering with letters warning of the dangers to public health. After the death of the workers, dozens of newspapers reported on the issue. The New York Times editorialized in 1924 that the deaths should not interfere with the production of more powerful fuel.
To settle the issue, the U.S. Public Health Service conducted a conference in 1925, and the sales of TEL were voluntarily suspended for one year to conduct a hazard assessment. The conference was initially expected to last for several days, but reportedly the conference decided that evaluating presentations on alternative anti-knock agents was not "its province", so it lasted a single day. Kettering and Midgley stated that no alternatives for anti-knocking were available, although private memos showed discussion of such agents. One commonly discussed agent was ethanol. The Public Health Service created a committee that reviewed a government-sponsored study of workers and an Ethyl lab test, and concluded that while leaded gasoline should not be banned, it should continue to be investigated. The low concentrations present in gasoline and exhaust were not perceived as immediately dangerous. A U.S. Surgeon General committee issued a report in 1926 that concluded there was no real evidence that the sale of TEL was hazardous to human health but urged further study. In the years that followed, research was heavily funded by the lead industry; in 1943, Randolph Byers found children with lead poisoning had behavior problems, but the Lead Industries Association threatened him with a lawsuit and the research ended.
In the late 1920s, Robert A. Kehoe of the University of Cincinnati was the Ethyl Corporation's chief medical consultant and one of the lead industry's staunchest advocates, who would not be discredited until decades later by Dr. Clair Patterson's work on human lead burdens (see below) and other studies. In 1928, Dr. Kehoe expressed the opinion that there was no basis for concluding that leaded fuels posed any health threat. He convinced the Surgeon General that the dose–response relationship of lead had "no effect" below a certain threshold. As the head of Kettering Laboratories for many years, Kehoe would become a chief promoter of the safety of TEL, an influence that did not begin to wane until about the early 1960s. But by the 1970s, the general opinion of the safety of TEL would change, and by 1976 the U.S. government would begin to require the phaseout of this product.
In the late 1940s and early 1950s, Clair Cameron Patterson accidentally discovered the pollution caused by TEL in the environment while determining the age of the Earth. As he attempted to measure lead content of very old rocks, and the time it took uranium to decay into lead, the readings were made inaccurate by lead in the environment that contaminated his samples. He was then forced to work in a cleanroom to keep his samples uncontaminated by environmental pollution of lead. After coming up with a fairly accurate estimate of the age of the Earth, he turned to investigating the lead contamination problem by examining ice cores from countries such as Greenland. He realized that the lead contamination in the environment dated from about the time that TEL became widely used as a fuel additive in gasoline. Being aware of the health dangers posed by lead and suspicious of the pollution caused by TEL, he became one of the earliest and most effective proponents of removing it from use.
In the 1960s, the first clinical works were published proving the toxicity of this compound in humans, e.g. by Mirosław Jan Stasik.
Modern findings
In the 1970s, Herbert Needleman found that higher lead levels in children were correlated with decreased school performance. Needleman was repeatedly accused of scientific misconduct by individuals within the lead industry, but he was eventually cleared by a scientific advisory council. Needleman also wrote the average US child's blood lead level was 13.7 μg/dL in 1976 and that Patterson believed that everyone was to some degree poisoned by TEL in gasoline.
In the U.S. in 1973, the United States Environmental Protection Agency issued regulations to reduce the lead content of leaded gasoline over a series of annual phases, which therefore came to be known as the "lead phasedown" program. EPA's rules were issued under section 211 of the Clean Air Act, as amended 1970. The Ethyl Corp challenged the EPA regulations in Federal court. Although the EPA's regulation was initially invalidated, the EPA won the case on appeal, so the TEL phasedown began to be implemented in 1976. Leaded gas was banned in vehicles with catalytic converters in 1975 due to damage of catalytic converters but it continued to be sold for vehicles without catalytic converters. Additional regulatory changes were made by EPA over the next decade (including adoption of a trading market in "lead credits" in 1982 that became the precursor of the Acid Rain Allowance Market, adopted in 1990 for SO2), but the decisive rule was issued in 1985. The EPA mandated that lead additive be reduced by 91 percent by the end of 1986. A 1994 study had indicated that the concentration of lead in the blood of the U.S. population had dropped 78% from 1976 to 1991. The U.S. phasedown regulations also were due in great part to studies conducted by Philip J. Landrigan.
In Europe, Professor Derek Bryce-Smith was among the first to highlight the potential dangers of TEL and became a leading campaigner for removal of lead additives from petrol.
From 1 January 1996, the U.S. Clean Air Act banned the sale of leaded fuel for use in on-road vehicles although that year the US EPA indicated that TEL could still be used in aircraft, racing cars, farm equipment, and marine engines. Thus, what had begun in the U.S. as a phasedown ultimately ended in a phase-out for on-road vehicle TEL. Similar bans in other countries have resulted in lowering levels of lead in people's bloodstreams.
Taking cue from the domestic programs, the U.S. Agency for International Development undertook an initiative to reduce tetraethyl lead use in other countries, notably its efforts in Egypt begun in 1995. In 1996, with the cooperation of the U.S. AID, Egypt took almost all of the lead out of its gasoline. The success in Egypt provided a model for AID efforts worldwide.
By 2000, the TEL industry had moved the major portion of their sales to developing countries whose governments they lobbied against phasing out leaded gasoline. Leaded gasoline was withdrawn entirely from the European Union market on 1 January 2000, although it had been banned earlier in most member states. Other countries also phased out TEL. India banned leaded petrol in March 2000.
By 2011, the United Nations announced that it had been successful in phasing out leaded gasoline worldwide. "Ridding the world of leaded petrol, with the United Nations leading the effort in developing countries, has resulted in $2.4 trillion in annual benefits, 1.2 million fewer premature deaths, higher overall intelligence and 58 million fewer crimes", the United Nations Environmental Programme said. The announcement was slightly premature, as a few countries still had leaded gasoline for sale as of 2017. On 30 August 2021 the United Nations Environment Programme announced that leaded gasoline had been eliminated. The final stocks of the product were used up in Algeria, which had continued to produce leaded gasoline until July 2021.
Effect on crime rates
Reduction in the average blood lead level is believed to have been a major cause for falling violent crime rates in the United States. A statistically significant correlation has been found between the usage rate of leaded gasoline and violent crime: the violent crime curve virtually tracks the lead exposure curve with a 22-year time lag. After the ban on TEL, blood lead levels in U.S. children dramatically decreased. Researchers including Amherst College economist Jessica Wolpaw Reyes, Department of Housing and Urban Development consultant Rick Nevin, and Howard Mielke of Tulane University say that declining exposure to lead is responsible for an up to 56% decline in crime from 1992 to 2002. Taking into consideration other factors that are believed to have increased crime rates over that period, Reyes found that the reduced exposure to lead led to an actual decline of 34% over that period.
Lingering issues over time
Although leaded gasoline has long since ended its history of regular use in U.S. transportation, it has left high concentrations of lead in the soil adjacent to roads that were heavily used prior to its phaseout. These contaminated materials present health dangers even when merely touched or when components of it get breathed in. Children, especially those in poverty inside of the U.S., are particularly at risk.
| Physical sciences | Organometallics | Chemistry |
3498698 | https://en.wikipedia.org/wiki/Caustic%20%28optics%29 | Caustic (optics) | In optics, a caustic or caustic network is the envelope of light rays which have been reflected or refracted by a curved surface or object, or the projection of that envelope of rays on another surface. The caustic is a curve or surface to which each of the light rays is tangent, defining a boundary of an envelope of rays as a curve of concentrated light. In some cases caustics can be seen as patches of light or their bright edges, shapes which often have cusp singularities.
Explanation
Concentration of light, especially sunlight, can burn. The word caustic, in fact, comes from the Greek καυστός, burnt, via the Latin causticus, burning.
A common situation where caustics are visible is when light shines on a drinking glass. The glass casts a shadow, but also produces a curved region of bright light. In ideal circumstances (including perfectly parallel rays, as if from a point source at infinity), a nephroid-shaped patch of light can be produced. Rippling caustics are commonly formed when light shines through waves on a body of water.
Another familiar caustic is the rainbow. Scattering of light by raindrops causes different wavelengths of light to be refracted into arcs of differing radius, producing the bow.
Computer graphics
In computer graphics, most modern rendering systems support caustics. Some of them even support volumetric caustics. This is accomplished by raytracing the possible paths of a light beam, accounting for the refraction and reflection. Photon mapping is one implementation of this. Volumetric caustics can also be achieved by volumetric path tracing. Some computer graphic systems work by "forward ray tracing" wherein photons are modeled as coming from a light source and bouncing around the environment according to rules. Caustics are formed in the regions where sufficient photons strike a surface causing it to be brighter than the average area in the scene. “Backward ray tracing” works in the reverse manner beginning at the surface and determining if there is a direct path to the light source. Some examples of 3D ray-traced caustics can be found here.
The focus of most computer graphics systems is aesthetics rather than physical accuracy. This is especially true when it comes to real-time graphics in computer games where generic pre-calculated textures are mostly used instead of physically correct calculations.
Caustic engineering
Caustic engineering describes the process of solving the inverse problem to computer graphics. That is, given a specific image, to determine a surface whose refracted or reflected light forms this image.
In the discrete version of this problem, the surface is divided into several micro-surfaces which are assumed smooth, i.e. the light reflected/refracted by each micro-surface forms a Gaussian caustic. Gaussian caustic means that each micro-surface obey Gaussian distribution. The position and orientation of each of the micro-surfaces are then obtained using a combination of Poisson integration and simulated annealing.
There have been many different approaches to address the continuous problem. One approach uses an idea from transportation theory called optimal transport to find a mapping between incoming light rays and the target surface. After obtaining such a mapping, the surface is optimized by adapting it iteratively using Snell's law of refraction.
Optimal-transport-based caustic pattern design
Basic principle
Controlling caustic pattern is rather a challenging problem as very minor changes of the surface will significantly affect the quality of the pattern since light ray directions might be interfered by other light rays as they intersect with and refract through the material. This will lead to a scattered, discontinuous pattern. To tackle this problem, optimal-transport-based is one of the existing proposed methods to control caustic pattern by redirecting light's directions as it propagates through the surface of a certain transparent material. This is done by solving an inverse optimization problem based on optimal transport. Given a reference image of an object/pattern, the target is to formulate the mathematical description of the material surface through which light refracts and converges to the similar pattern of the reference image. This is done by rearranging/recomputing the initial light intensity until the minimum of the optimization problem is reached.
Manufacturing
Once the caustic pattern has been designed computationally, the processed data will be then sent to the manufacturing stage to get the final product. The most common approach is subtractive manufacturing (machining).
Various materials can be used depending on the desired quality, the effort it takes to manufacture, and the available manufacturing method.
Common refractive materials: acrylic, polycarbonate, polyethylene, glass, diamond
Common reflective materials: steel, iron, aluminum, gold, silver, titanium, nickel
Caustic pattern design has many real-world applications, for example in:
Luminaires
Jewelry
Architecture
Decorative glass production
| Physical sciences | Optics | Physics |
3503207 | https://en.wikipedia.org/wiki/Multi-core%20processor | Multi-core processor | A multi-core processor (MCP) is a microprocessor on a single integrated circuit (IC) with two or more separate central processing units (CPUs), called cores to emphasize their multiplicity (for example, dual-core or quad-core). Each core reads and executes program instructions, specifically ordinary CPU instructions (such as add, move data, and branch). However, the MCP can run instructions on separate cores at the same time, increasing overall speed for programs that support multithreading or other parallel computing techniques. Manufacturers typically integrate the cores onto a single IC die, known as a chip multiprocessor (CMP), or onto multiple dies in a single chip package. As of 2024, the microprocessors used in almost all new personal computers are multi-core.
A multi-core processor implements multiprocessing in a single physical package. Designers may couple cores in a multi-core device tightly or loosely. For example, cores may or may not share caches, and they may implement message passing or shared-memory inter-core communication methods. Common network topologies used to interconnect cores include bus, ring, two-dimensional mesh, and crossbar. Homogeneous multi-core systems include only identical cores; heterogeneous multi-core systems have cores that are not identical (e.g. big.LITTLE have heterogeneous cores that share the same instruction set, while AMD Accelerated Processing Units have cores that do not share the same instruction set). Just as with single-processor systems, cores in multi-core systems may implement architectures such as VLIW, superscalar, vector, or multithreading.
Multi-core processors are widely used across many application domains, including general-purpose, embedded, network, digital signal processing (DSP), and graphics (GPU). Core count goes up to even dozens, and for specialized chips over 10,000, and in supercomputers (i.e. clusters of chips) the count can go over 10 million (and in one case up to 20 million processing elements total in addition to host processors).
The improvement in performance gained by the use of a multi-core processor depends very much on the software algorithms used and their implementation. In particular, possible gains are limited by the fraction of the software that can run in parallel simultaneously on multiple cores; this effect is described by Amdahl's law. In the best case, so-called embarrassingly parallel problems may realize speedup factors near the number of cores, or even more if the problem is split up enough to fit within each core's cache(s), avoiding use of much slower main-system memory. Most applications, however, are not accelerated as much unless programmers invest effort in refactoring.
The parallelization of software is a significant ongoing topic of research. Cointegration of multiprocessor applications provides flexibility in network architecture design. Adaptability within parallel models is an additional feature of systems utilizing these protocols.
In the consumer market, dual-core processors (that is, microprocessors with two units) started becoming commonplace on personal computers in the late 2000s. Quad-core processors were also being adopted in that era for higher-end systems before becoming standard. In the late 2010s, hexa-core (six cores) started entering the mainstream and since the early 2020s has overtaken quad-core in many spaces.
Terminology
The terms multi-core and dual-core most commonly refer to some sort of central processing unit (CPU), but are sometimes also applied to digital signal processors (DSP) and system on a chip (SoC). The terms are generally used only to refer to multi-core microprocessors that are manufactured on the same integrated circuit die; separate microprocessor dies in the same package are generally referred to by another name, such as multi-chip module. This article uses the terms "multi-core" and "dual-core" for CPUs manufactured on the same integrated circuit, unless otherwise noted.
In contrast to multi-core systems, the term multi-CPU refers to multiple physically separate processing-units (which often contain special circuitry to facilitate communication between each other).
The terms many-core and massively multi-core are sometimes used to describe multi-core architectures with an especially high number of cores (tens to thousands).
Some systems use many soft microprocessor cores placed on a single FPGA. Each "core" can be considered a "semiconductor intellectual property core" as well as a CPU core.
Development
While manufacturing technology improves, reducing the size of individual gates, physical limits of semiconductor-based microelectronics have become a major design concern. These physical limitations can cause significant heat dissipation and data synchronization problems. Various other methods are used to improve CPU performance. Some instruction-level parallelism (ILP) methods such as superscalar pipelining are suitable for many applications, but are inefficient for others that contain difficult-to-predict code. Many applications are better suited to thread-level parallelism (TLP) methods, and multiple independent CPUs are commonly used to increase a system's overall TLP. A combination of increased available space (due to refined manufacturing processes) and the demand for increased TLP led to the development of multi-core CPUs.
Commercial incentives
Several business motives drive the development of multi-core architectures. For decades, it was possible to improve performance of a CPU by shrinking the area of the integrated circuit (IC), which reduced the cost per device on the IC. Alternatively, for the same circuit area, more transistors could be used in the design, which increased functionality, especially for complex instruction set computing (CISC) architectures. Clock rates also increased by orders of magnitude in the decades of the late 20th century, from several megahertz in the 1980s to several gigahertz in the early 2000s.
As the rate of clock speed improvements slowed, increased use of parallel computing in the form of multi-core processors has been pursued to improve overall processing performance. Multiple cores were used on the same CPU chip, which could then lead to better sales of CPU chips with two or more cores. For example, Intel has produced a 48-core processor for research in cloud computing; each core has an x86 architecture.
Technical factors
Since computer manufacturers have long implemented symmetric multiprocessing (SMP) designs using discrete CPUs, the issues regarding implementing multi-core processor architecture and supporting it with software are well known.
Additionally:
Using a proven processing-core design without architectural changes reduces design risk significantly.
For general-purpose processors, much of the motivation for multi-core processors comes from greatly diminished gains in processor performance from increasing the operating frequency. This is due to three primary factors:
The memory wall; the increasing gap between processor and memory speeds. This, in effect, pushes for cache sizes to be larger in order to mask the latency of memory. This helps only to the extent that memory bandwidth is not the bottleneck in performance.
The ILP wall; the increasing difficulty of finding enough parallelism in a single instruction stream to keep a high-performance single-core processor busy.
The power wall; the trend of consuming exponentially increasing power (and thus also generating exponentially increasing heat) with each factorial increase of operating frequency. This increase can be mitigated by "shrinking" the processor by using smaller traces for the same logic. The power wall poses manufacturing, system design and deployment problems that have not been justified in the face of the diminished gains in performance due to the memory wall and ILP wall.
In order to continue delivering regular performance improvements for general-purpose processors, manufacturers such as Intel and AMD have turned to multi-core designs, sacrificing lower manufacturing-costs for higher performance in some applications and systems. Multi-core architectures are being developed, but so are the alternatives. An especially strong contender for established markets is the further integration of peripheral functions into the chip.
Advantages
The proximity of multiple CPU cores on the same die allows the cache coherency circuitry to operate at a much higher clock rate than what is possible if the signals have to travel off-chip. Combining equivalent CPUs on a single die significantly improves the performance of cache snoop (alternative: Bus snooping) operations. Put simply, this means that signals between different CPUs travel shorter distances, and therefore those signals degrade less. These higher-quality signals allow more data to be sent in a given time period, since individual signals can be shorter and do not need to be repeated as often.
Assuming that the die can physically fit into the package, multi-core CPU designs require much less printed circuit board (PCB) space than do multi-chip SMP designs. Also, a dual-core processor uses slightly less power than two coupled single-core processors, principally because of the decreased power required to drive signals external to the chip. Furthermore, the cores share some circuitry, like the L2 cache and the interface to the front-side bus (FSB). In terms of competing technologies for the available silicon die area, multi-core design can make use of proven CPU core library designs and produce a product with lower risk of design error than devising a new wider-core design. Also, adding more cache suffers from diminishing returns.
Multi-core chips also allow higher performance at lower energy. This can be a big factor in mobile devices that operate on batteries. Since each core in a multi-core CPU is generally more energy-efficient, the chip becomes more efficient than having a single large monolithic core. This allows higher performance with less energy. A challenge in this, however, is the additional overhead of writing parallel code.
Disadvantages
Maximizing the usage of the computing resources provided by multi-core processors requires adjustments both to the operating system (OS) support and to existing application software. Also, the ability of multi-core processors to increase application performance depends on the use of multiple threads within applications.
Integration of a multi-core chip can lower the chip production yields. They are also more difficult to manage thermally than lower-density single-core designs. Intel has partially countered this first problem by creating its quad-core designs by combining two dual-core ones on a single die with a unified cache, hence any two working dual-core dies can be used, as opposed to producing four cores on a single die and requiring all four to work to produce a quad-core CPU. From an architectural point of view, ultimately, single CPU designs may make better use of the silicon surface area than multiprocessing cores, so a development commitment to this architecture may carry the risk of obsolescence. Finally, raw processing power is not the only constraint on system performance. Two processing cores sharing the same system bus and memory bandwidth limits the real-world performance advantage.
Hardware
Trends
The trend in processor development has been towards an ever-increasing number of cores, as processors with hundreds or even thousands of cores become theoretically possible. In addition, multi-core chips mixed with simultaneous multithreading, memory-on-chip, and special-purpose "heterogeneous" (or asymmetric) cores promise further performance and efficiency gains, especially in processing multimedia, recognition and networking applications. For example, a big.LITTLE core includes a high-performance core (called 'big') and a low-power core (called 'LITTLE'). There is also a trend towards improving energy-efficiency by focusing on performance-per-watt with advanced fine-grain or ultra fine-grain power management and dynamic voltage and frequency scaling (i.e. laptop computers and portable media players).
Chips designed from the outset for a large number of cores (rather than having evolved from single core designs) are sometimes referred to as manycore designs, emphasising qualitative differences.
Architecture
The composition and balance of the cores in multi-core architecture show great variety. Some architectures use one core design repeated consistently ("homogeneous"), while others use a mixture of different cores, each optimized for a different, "heterogeneous" role.
How multiple cores are implemented and integrated significantly affects both the developer's programming skills and the consumer's expectations of apps and interactivity versus the device. A device advertised as being octa-core will only have independent cores if advertised as True Octa-core, or similar styling, as opposed to being merely two sets of quad-cores each with fixed clock speeds.
The article "CPU designers debate multi-core future" by Rick Merritt, EE Times 2008, includes these comments:
Software effects
An outdated version of an anti-virus application may create a new thread for a scan process, while its GUI thread waits for commands from the user (e.g. cancel the scan). In such cases, a multi-core architecture is of little benefit for the application itself due to the single thread doing all the heavy lifting and the inability to balance the work evenly across multiple cores. Programming truly multithreaded code often requires complex co-ordination of threads and can easily introduce subtle and difficult-to-find bugs due to the interweaving of processing on data shared between threads (see thread-safety). Consequently, such code is much more difficult to debug than single-threaded code when it breaks. There has been a perceived lack of motivation for writing consumer-level threaded applications because of the relative rarity of consumer-level demand for maximum use of computer hardware. Also, serial tasks like decoding the entropy encoding algorithms used in video codecs are impossible to parallelize because each result generated is used to help create the next result of the entropy decoding algorithm.
Given the increasing emphasis on multi-core chip design, stemming from the grave thermal and power consumption problems posed by any further significant increase in processor clock speeds, the extent to which software can be multithreaded to take advantage of these new chips is likely to be the single greatest constraint on computer performance in the future. If developers are unable to design software to fully exploit the resources provided by multiple cores, then they will ultimately reach an insurmountable performance ceiling.
The telecommunications market had been one of the first that needed a new design of parallel datapath packet processing because there was a very quick adoption of these multiple-core processors for the datapath and the control plane. These MPUs are going to replace the traditional Network Processors that were based on proprietary microcode or picocode.
Parallel programming techniques can benefit from multiple cores directly. Some existing parallel programming models such as Cilk Plus, OpenMP, OpenHMPP, FastFlow, Skandium, MPI, and Erlang can be used on multi-core platforms. Intel introduced a new abstraction for C++ parallelism called TBB. Other research efforts include the Codeplay Sieve System, Cray's Chapel, Sun's Fortress, and IBM's X10.
Multi-core processing has also affected the ability of modern computational software development. Developers programming in newer languages might find that their modern languages do not support multi-core functionality. This then requires the use of numerical libraries to access code written in languages like C and Fortran, which perform math computations faster than newer languages like C#. Intel's MKL and AMD's ACML are written in these native languages and take advantage of multi-core processing. Balancing the application workload across processors can be problematic, especially if they have different performance characteristics. There are different conceptual models to deal with the problem, for example using a coordination language and program building blocks (programming libraries or higher-order functions). Each block can have a different native implementation for each processor type. Users simply program using these abstractions and an intelligent compiler chooses the best implementation based on the context.
Managing concurrency acquires a central role in developing parallel applications. The basic steps in designing parallel applications are:
Partitioning The partitioning stage of a design is intended to expose opportunities for parallel execution. Hence, the focus is on defining a large number of small tasks in order to yield what is termed a fine-grained decomposition of a problem.
Communication The tasks generated by a partition are intended to execute concurrently but cannot, in general, execute independently. The computation to be performed in one task will typically require data associated with another task. Data must then be transferred between tasks so as to allow computation to proceed. This information flow is specified in the communication phase of a design.
Agglomeration In the third stage, development moves from the abstract toward the concrete. Developers revisit decisions made in the partitioning and communication phases with a view to obtaining an algorithm that will execute efficiently on some class of parallel computer. In particular, developers consider whether it is useful to combine, or agglomerate, tasks identified by the partitioning phase, so as to provide a smaller number of tasks, each of greater size. They also determine whether it is worthwhile to replicate data and computation.
Mapping In the fourth and final stage of the design of parallel algorithms, the developers specify where each task is to execute. This mapping problem does not arise on uniprocessors or on shared-memory computers that provide automatic task scheduling.
On the other hand, on the server side, multi-core processors are ideal because they allow many users to connect to a site simultaneously and have independent threads of execution. This allows for Web servers and application servers that have much better throughput.
Licensing
Vendors may license some software "per processor". This can give rise to ambiguity, because a "processor" may consist either of a single core or of a combination of cores.
Initially, for some of its enterprise software, Microsoft continued to use a per-socket licensing system. However, for some software such as BizTalk Server 2013, SQL Server 2014, and Windows Server 2016, Microsoft has shifted to per-core licensing.
Oracle Corporation counts an AMD X2 or an Intel dual-core CPU as a single processor but uses other metrics for other types, especially for processors with more than two cores.
Embedded applications
Embedded computing operates in an area of processor technology distinct from that of "mainstream" PCs. The same technological drives towards multi-core apply here too. Indeed, in many cases the application is a "natural" fit for multi-core technologies, if the task can easily be partitioned between the different processors.
In addition, embedded software is typically developed for a specific hardware release, making issues of software portability, legacy code or supporting independent developers less critical than is the case for PC or enterprise computing. As a result, it is easier for developers to adopt new technologies and as a result there is a greater variety of multi-core processing architectures and suppliers.
Network processors
, multi-core network processors have become mainstream, with companies such as Freescale Semiconductor, Cavium Networks, Wintegra and Broadcom all manufacturing products with eight processors. For the system developer, a key challenge is how to exploit all the cores in these devices to achieve maximum networking performance at the system level, despite the performance limitations inherent in a symmetric multiprocessing (SMP) operating system. Companies such as 6WIND provide portable packet processing software designed so that the networking data plane runs in a fast path environment outside the operating system of the network device.
Digital signal processing
In digital signal processing the same trend applies: Texas Instruments has the three-core TMS320C6488 and four-core TMS320C5441, Freescale the four-core MSC8144 and six-core MSC8156 (and both have stated they are working on eight-core successors). Newer entries include the Storm-1 family from Stream Processors, Inc with 40 and 80 general purpose ALUs per chip, all programmable in C as a SIMD engine and Picochip with 300 processors on a single die, focused on communication applications.
Heterogeneous systems
In heterogeneous computing, where a system uses more than one kind of processor or cores, multi-core solutions are becoming more common: Xilinx Zynq UltraScale+ MPSoC has a quad-core ARM Cortex-A53 and dual-core ARM Cortex-R5. Software solutions such as OpenAMP are being used to help with inter-processor communication.
Mobile devices may use the ARM big.LITTLE architecture.
Hardware examples
Commercial
Adapteva Epiphany, a many-core processor architecture which allows up to 4096 processors on-chip, although only a 16-core version has been commercially produced.
Aeroflex Gaisler LEON3, a multi-core SPARC that also exists in a fault-tolerant version.
Ageia PhysX, a multi-core physics processing unit.
Ambric Am2045, a 336-core massively parallel processor array (MPPA)
AMD
A-Series, dual-, triple-, and quad-core of Accelerated Processor Units (APU).
Athlon 64 FX and Athlon 64 X2 single- and dual-core desktop processors.
Athlon II, dual-, triple-, and quad-core desktop processors.
FX-Series, quad-, 6-, and 8-core desktop processors.
Opteron, single-, dual-, quad-, 6-, 8-, 12-, and 16-core server/workstation processors.
Phenom, dual-, triple-, and quad-core processors.
Phenom II, dual-, triple-, quad-, and 6-core desktop processors.
Sempron, single-, dual-, and quad-core entry level processors.
Turion, single- and dual-core laptop processors.
Ryzen, dual-, quad-, 6-, 8-, 12-, 16-, 24-, 32-, and 64-core desktop, mobile, and embedded platform processors.
Epyc, quad-, 8-, 12-, 16-, 24-, 32-, and 64-core server and embedded processors.
Radeon and FireStream GPU/GPGPU.
Analog Devices Blackfin BF561, a symmetrical dual-core processor
ARM MPCore is a fully synthesizable multi-core container for ARM11 MPCore and ARM Cortex-A9 MPCore processor cores, intended for high-performance embedded and entertainment applications.
ASOCS ModemX, up to 128 cores, wireless applications.
Azul Systems
Vega 1, a 24-core processor, released in 2005.
Vega 2, a 48-core processor, released in 2006.
Vega 3, a 54-core processor, released in 2008.
Broadcom
SiByte SB1250, SB1255, SB1455
BCM2836, BCM2837, BCM2710 and BCM2711 quad-core ARM SoC (designed for different Raspberry Pi models)
Cadence Design Systems Tensilica Xtensa LX6, available in a dual-core configuration in Espressif Systems's ESP32
ClearSpeed
CSX700, 192-core processor, released in 2008 (32/64-bit floating point; Integer ALU).
Cradle Technologies CT3400 and CT3600, both multi-core DSPs.
Cavium Networks Octeon, a 32-core MIPS MPU.
Coherent Logix hx3100 Processor, a 100-core DSP/GPP processor.
Freescale Semiconductor QorIQ series processors, up to 8 cores, Power ISA MPU.
Hewlett-Packard PA-8800 and PA-8900, dual core PA-RISC processors.
IBM
POWER4, a dual-core PowerPC processor, released in 2001.
POWER5, a dual-core PowerPC processor, released in 2004.
POWER6, a dual-core PowerPC processor, released in 2007.
POWER7, a 4, 6 and 8-core PowerPC processor, released in 2010.
POWER8, a 12-core PowerPC processor, released in 2013.
POWER9, a 12 or 24-core PowerPC processor, released in 2017.
Power10, a 15 or 30-core PowerPC processor, released in 2021.
PowerPC 970MP, a dual-core PowerPC processor, used in the Apple Power Mac G5.
Xenon, a triple-core, SMT-capable, PowerPC microprocessor used in the Microsoft Xbox 360 game console.
z10, a quad-core z/Architecture processor, released in 2008.
z196, a quad-core z/Architecture processor, released in 2010.
zEC12, a six-core z/Architecture processor, released in 2012.
z13, an eight-core z/Architecture processor, released in 2015.
z14, a ten-core z/Architecture processor, released in 2017.
z15, a twelve-core z/Architecture processor, released in 2019.
Telum, an eight-core z/Architecture processor, released in 2021.
Infineon
AURIX
Danube, a dual-core, MIPS-based, home gateway processor.
Intel
Atom, single, dual-core, quad-core, 8-, 12-, and 16-core processors for netbooks, nettops, embedded applications, and mobile internet devices (MIDs).
Atom SoC (system on a chip), single-core, dual-core, and quad-core processors for smartphones and tablets.
Celeron, the first dual-core (and, later, quad-core) processor for the budget/entry-level market.
Core Duo, a dual-core processor.
Core 2 Duo, a dual-core processor.
Core 2 Quad, 2 dual-core dies packaged in a multi-chip module.
Core i3, Core i5, Core i7 and Core i9, a family of dual-, quad-, 6-, 8-, 10-, 12-, 14-, 16-, and 18-core processors, and the successor of the Core 2 Duo and the Core 2 Quad.
Itanium, single, dual-core, quad-core, and 8-core processors.
Pentium, single, dual-core, and quad-core processors for the entry-level market.
Teraflops Research Chip (Polaris), a 3.16 GHz, 80-core processor prototype, which the company originally stated would be released by 2011.
Xeon dual-, quad-, 6-, 8-, 10-, 12-, 14-, 15-, 16-, 18-, 20-, 22-, 24-, 26-, 28-, 32-, 48-, and 56-core processors.
Xeon Phi 57-, 60-, 61-, 64-, 68-, and 72-core processors.
IntellaSys
SEAforth 40C18, a 40-core processor.
SEAforth24, a 24-core processor designed by Charles H. Moore.
Kalray
MPPA-256, 256-core processor, released 2012 (256 usable VLIW cores, Network-on-Chip (NoC), 32/64-bit IEEE 754 compliant FPU)
NetLogic Microsystems
XLP, a 32-core, quad-threaded MIPS64 processor.
XLR, an eight-core, quad-threaded MIPS64 processor.
XLS, an eight-core, quad-threaded MIPS64 processor.
Nvidia
RTX 3090 (10496 CUDA cores, GPGPU cores; plus other more specialized cores).
Parallax Propeller P8X32, an eight-core microcontroller.
picoChip PC200 series 200–300 cores per device for DSP & wireless.
Plurality HAL series tightly coupled 16-256 cores, L1 shared memory, hardware synchronized processor.
Rapport Kilocore KC256, a 257-core microcontroller with a PowerPC core and 256 8-bit "processing elements".
Raspberry Pi Ltd. RP2040, a dual ARM Cortex-M0+ microcontroller
SiCortex "SiCortex node" has six MIPS64 cores on a single chip.
SiFive
U74 includes 4 cores
Sony/IBM/Toshiba's Cell processor, a nine-core processor with one general purpose PowerPC core and eight specialized SPUs (Synergistic Processing Unit) optimized for vector operations used in the Sony PlayStation 3.
Sun Microsystems
MAJC 5200, two-core VLIW processor.
UltraSPARC IV and UltraSPARC IV+, dual-core processors.
UltraSPARC T1, an eight-core, 32-thread processor.
UltraSPARC T2, an eight-core, 64-concurrent-thread processor.
UltraSPARC T3, a sixteen-core, 128-concurrent-thread processor.
SPARC T4, an eight-core, 64-concurrent-thread processor.
SPARC T5, a sixteen-core, 128-concurrent-thread processor.
Sunway
Sunway SW26010, a 260-core processor used in the Sunway TaihuLight.
Texas Instruments
TMS320C80 MVP, a five-core multimedia video processor.
TMS320TMS320C66, 2-, 4-, 8-core DSP.
Tilera
TILE64, a 64-core 32-bit processor.
TILE-Gx, a 72-core 64-bit processor.
XMOS Software Defined Silicon quad-core XS1-G4.
Free
OpenSPARC
Academic
Stanford, 4-core Hydra processor
MIT, 16-core RAW processor
University of California, Davis, Asynchronous array of simple processors (AsAP)
36-core 610 MHz AsAP
167-core 1.2 GHz AsAP2
University of Washington, Wavescalar processor
University of Texas, Austin, TRIPS processor
Linköping University, Sweden, ePUMA processor
UC Davis, Kilocore, a 1000 core 1.78 GHz processor on a 32 nm IBM process
Benchmarks
The research and development of multicore processors often compares many options, and benchmarks are developed to help such evaluations. Existing benchmarks include SPLASH-2, PARSEC, and COSMIC for heterogeneous systems.
| Technology | Computer hardware | null |
23716451 | https://en.wikipedia.org/wiki/Generalized%20Petersen%20graph | Generalized Petersen graph | In graph theory, the generalized Petersen graphs are a family of cubic graphs formed by connecting the vertices of a regular polygon to the corresponding vertices of a star polygon. They include the Petersen graph and generalize one of the ways of constructing the Petersen graph. The generalized Petersen graph family was introduced in 1950 by H. S. M. Coxeter and was given its name in 1969 by Mark Watkins.
Definition and notation
In Watkins' notation, G(n, k) is a graph with vertex set
and edge set
where subscripts are to be read modulo n and k < n/2. Some authors use the notation GPG(n, k). Coxeter's notation for the same graph would be {n} + {n/k}, a combination of the Schläfli symbols for the regular n-gon and star polygon from which the graph is formed. The Petersen graph itself is G(5, 2) or {5} + {5/2}.
Any generalized Petersen graph can also be constructed from a voltage graph with two vertices, two self-loops, and one other edge.
Examples
Among the generalized Petersen graphs are the n-prism G(n, 1), the Dürer graph G(6, 2), the Möbius-Kantor graph G(8, 3), the dodecahedron G(10, 2), the Desargues graph G(10, 3) and the Nauru graph G(12, 5).
Four generalized Petersen graphs – the 3-prism, the 5-prism, the Dürer graph, and G(7, 2) – are among the seven graphs that are cubic, 3-vertex-connected, and well-covered (meaning that all of their maximal independent sets have equal size).
Properties
This family of graphs possesses a number of interesting properties. For example:
G(n, k) is vertex-transitive (meaning that it has symmetries that take any vertex to any other vertex) if and only if (n, k) = (10, 2) or k2 ≡ ±1 (mod n).
G(n, k) is edge-transitive (having symmetries that take any edge to any other edge) only in the following seven cases: (n, k) = (4, 1), (5, 2), (8, 3), (10, 2), (10, 3), (12, 5), (24, 5). These seven graphs are therefore the only symmetric generalized Petersen graphs.
G(n, k) is bipartite if and only if n is even and k is odd.
G(n, k) is a Cayley graph if and only if k2 ≡ 1 (mod n).
G(n, k) is hypohamiltonian when n is congruent to 5 modulo 6 and k = 2, n − 2, or (n ± 1)/2 (these four choices of k lead to isomorphic graphs). It is also non-Hamiltonian when n is divisible by 4, at least equal to 8, and k = n/2. In all other cases it has a Hamiltonian cycle. When n is congruent to 3 modulo 6 G(n, 2) has exactly three Hamiltonian cycles. For G(n, 2), the number of Hamiltonian cycles can be computed by a formula that depends on the congruence class of n modulo 6 and involves the Fibonacci numbers.
Every generalized Petersen graph is a unit distance graph.
Isomorphisms
G(n, k) is isomorphic to G(n, l) if and only if k=l or kl ≡ ±1 (mod n).
Girth
The girth of G(n, k) is at least 3 and at most 8, in particular:
A table with exact girth values:
Chromatic number and chromatic index
Generalized Petersen graphs are regular graphs of degree three, so according to Brooks' theorem their chromatic number can only be two or three. More exactly:
Where denotes the logical AND, while the logical OR. Here, denotes divisibility, and denotes its negation. For example, the chromatic number of is 3.
The Petersen graph, being a snark, has a chromatic index of 4: its edges require four colors. All other generalized Petersen graphs have chromatic index 3. These are the only possibilities, by Vizing's theorem.
The generalized Petersen graph G(9, 2) is one of the few graphs known to have only one 3-edge-coloring.
The Petersen graph itself is the only generalized Petersen graph that is not 3-edge-colorable.
Perfect Colorings
All admissible matrices of all perfect 2-colorings of the graphs G(n, 2) and G(n, 3) are enumerated.
| Mathematics | Graph theory | null |
20727645 | https://en.wikipedia.org/wiki/Somatic%20fusion | Somatic fusion | Somatic fusion, also called protoplast fusion, is a type of genetic modification in plants by which two distinct species of plants are fused together to form a new hybrid plant with the characteristics of both, a somatic hybrid. Hybrids have been produced either between different varieties of the same species (e.g. between non-flowering potato plants and flowering potato plants) or between two different species (e.g. between wheat Triticum and rye Secale to produce Triticale).
Uses of somatic fusion include developing plants resistant to disease, such as making potato plants resistant to potato leaf roll disease. Through somatic fusion, the crop potato plant Solanum tuberosum – the yield of which is severely reduced by a viral disease transmitted on by the aphid vector – is fused with the wild, non-tuber-bearing potato Solanum brevidens, which is resistant to the disease. The resulting hybrid has the chromosomes of both plants and is thus similar to polyploid plants.
Somatic hybridization was first introduced by Carlson et al. in Nicotiana glauca.
Process for plant cells
The somatic fusion process occurs in four steps:
The removal of the cell wall of one cell of each type of plant using cellulase enzyme to produce a somatic cell called a protoplast
The cells are then fused using electric shock (electrofusion) or chemical treatment to join the cells and fuse together the nuclei. The resulting fused nucleus is called heterokaryon.
The formation of the cell wall is then induced using hormones
The cells are then grown into calluses which then are further grown to plantlets and finally to a full plant, known as a somatic hybrid.
The procedure for seed plants describe above, fusion of moss protoplasts can be initiated without electric shock but by the use of polyethylene glycol (PEG). Further, moss protoplasts do not need phytohormones for regeneration, and they do not form a callus. Instead, regenerating moss protoplasts behave like germinating moss spores. Of further note sodium nitrate and calcium ion at high pH can be used, although results are variable depending on the organism.
Applications of hybrid cells
Somatic cells of different types can be fused to obtain hybrid cells. Hybrid cells are useful in a variety of ways, e.g.,
(i) to study the control of cell division and gene expression,
(ii) to investigate malignant transformations,
(iii) to obtain viral replication,
(iv) for gene or chromosome mapping and for
(v) production of monoclonal antibodies by producing hybridoma (hybrid cells between an immortalised cell and an antibody producing lymphocyte), etc.
Chromosome mapping through somatic cell hybridization is essentially based on fusion of human and mouse somatic cells. Generally, human fibrocytes or leucocytes are fused with mouse continuous cell lines.
When human and mouse cells (or cells of any two mammalian species or of the same species) are mixed, spontaneous cell fusion occurs at a very low rate (10-6). Cell fusion is enhanced 100 to 1000 times by the addition of ultraviolet inactivated Sendai (parainfluenza) virus or polyethylene glycol (PEG).
These agents adhere to the plasma membranes of cells and alter their properties in such a way that facilitates their fusion. Fusion of two cells produces a heterokaryon, i.e., a single hybrid cell with two nuclei, one from each of the cells entering fusion. Subsequently, the two nuclei also fuse to yield a hybrid cell with a single nucleus.
A generalized scheme for somatic cell hybridization may be described as follows. Appropriate human and mouse cells are selected and mixed together in the presence of inactivated Sendai virus or PEG to promote cell fusion. After a period of time, the cells (a mixture of man, mouse and 'hybrid' cells) are plated on a selective medium, e.g., HAT medium, which allows the multiplication of hybrid cells only.
Several clones (each derived from a single hybrid cell) of the hybrid cells are thus isolated and subjected to both cytogenetic and appropriate biochemical analyses for the detection of enzyme/ protein/trait under investigation. An attempt is now made to correlate the presence and absence of the trait with the presence and absence of a human chromosome in the hybrid clones.
If there is a perfect correlation between the presence and absence of a human chromosome and that of a trait in the hybrid clones, the gene governing the trait is taken to be located in the concerned chromosome.
The HAT medium is one of the several selective media used for the selection of hybrid cells. This medium is supplemented with hypoxanthine, aminopterin and thymidine, hence the name HAT medium. Antimetabolite aminopterin blocks the cellular biosynthesis of purines and pyrimidines from simple sugars and amino acids.
However, normal human and mouse cells can still multiply as they can utilize hypoxanthine and thymidine present in the medium through a salvage pathway, which ordinarily recycles the purines and pyrimidines produced from degradation of nucleic acids.
Hypoxanthine is converted into guanine by the enzyme hypoxanthine-guanine phosphoribosyltransferase (HGPRT), while thymidine is phosphorylated by thymidine kinase (TK); both HGPRT and TK are enzymes of the salvage pathway.
On a HAT medium, only those cells that have active HGPRT (HGPRT+) and TK (TK+) enzymes can proliferate, while those deficient in these enzymes (HGPRr- and/or TK-) can not divide (since they cannot produce purines and pyrimidines due to the aminopterin present in the HAT medium).
For using HAT medium as a selective agent, human cells used for fusion must be deficient for either the enzyme HGPRT or TK, while mouse cells must be deficient for the other enzyme of this pair. Thus, one may fuse HGPRT deficient human cells (designated as TK+ HGPRr-) with TK deficient mouse cells (denoted as TK- HGPRT+).
Their fusion products (hybrid cells) will be TK+ (due to the human gene) and HGPRT+ (due to the mouse gene) and will multiply on the HAT medium, while the man and mouse cells will fail to do so. Experiments with other selective media can be planned in a similar fashion.
Characteristics of somatic hybridization and cybridization
Somatic cell fusion appears to be the only means through which two different parental genomes can be recombined among plants that cannot reproduce sexually (asexual or sterile).
Protoplasts of sexually sterile (haploid, triploid, and aneuploid) plants can be fused to produce fertile diploids and polyploids.
Somatic cell fusion overcomes sexual incompatibility barriers. In some cases somatic hybrids between two incompatible plants have also found application in industry or agriculture.
Somatic cell fusion is useful in the study of cytoplasmic genes and their activities and this information can be applied in plant breeding experiments.
Inter-specific and inter-generic fusion achievements
Note: The table only lists a few examples, there are many more crosses. The possibilities of this technology are great; however, not all species are easily put into protoplast culture.
| Technology | Biotechnology | null |
20733467 | https://en.wikipedia.org/wiki/Bowhead%20whale | Bowhead whale | The bowhead whale (Balaena mysticetus), sometimes called the Greenland right whale, Arctic whale, and polar whale, is a species of baleen whale belonging to the family Balaenidae and is the only living representative of the genus Balaena. It is the only baleen whale endemic to the Arctic and subarctic waters, and is named after its characteristic massive triangular skull, which it uses to break through Arctic ice.
Bowheads have the largest mouth of any animal representing almost one-third of the length of the body, the longest baleen plates with a maximum length of , and may be the longest-lived mammals, with the ability to reach an age of more than 200 years.
The bowhead was an early whaling target. Their population was severely reduced before a 1966 moratorium was passed to protect the species. Of the five stocks of bowhead populations, three are listed as "endangered", one as "vulnerable", and one as "lower risk, conservation dependent" according to the IUCN Red List. The global population is assessed as of least concern.
Taxonomy
Carl Linnaeus named this species in the tenth edition of his Systema Naturae (1758). It was seemingly identical to its relatives in the North Atlantic, North Pacific, and Southern Oceans, and as such they were all thought to be a single species, collectively known as the "right whale", and given the binomial name Balaena mysticetus.
Today, the bowhead whale occupies a monotypic genus, separate from the right whales, as proposed by the work of John Edward Gray in 1821.
For the next 180 years, the family Balaenidae was the subject of great taxonometric debate. Authorities have repeatedly recategorized the three populations of right whale plus the bowhead whale, as one, two, three or four species, either in a single genus or in two separate genera. Eventually, it was recognized that bowheads and right whales were different, but there was still no strong consensus as to whether they shared a single genus or two. As recently as 1998, Dale Rice listed just two species – B. glacialis (the right whales) and B. mysticetus (the bowheads) – in his comprehensive and otherwise authoritative classification.
Studies in the 2000s finally provided clear evidence that the three living right whale species comprise a phylogenetic lineage, distinct from the bowhead, and that the bowhead and the right whales are rightly classified into two separate genera.
The right whales were thus confirmed to be in a separate genus, Eubalaena. The relationship is shown in the cladogram below:
The earlier fossil record shows no related cetacean after Morenocetus, found in a South American deposit dating back 23 million years.
An unknown species of right whale, the so-called "Swedenborg whale", which was proposed by Emanuel Swedenborg in the 18th century, was once thought to be a North Atlantic right whale. Based on later DNA analysis, those fossil bones claimed to be from Swedenborg whales were confirmed to be from bowhead whales.
Description
The bowhead whale is among the largest baleen whale species and is distinguished by its round body with an exceptionally curved rostrum, a large head, and long, dark baleen plates. Relative to its size, the bowhead whale has the largest head of any cetacean, measuring nearly 40% of the total body length. Two blowholes are situated atop its head, and help propell a stream of water up to in the air. The lower lips encompasses the baleen racks and resembles a curved circular shape when viewed from the side. It also has wide, trigonal flukes and fairly large, oar-shaped flippers. The skin is mostly black with white patches around the flukes, tail, eyes, and chin. These patches develop throughout life, with the exception of the patch surrounding the chin, which is usually visible in newborn calves, and increases in size at the same rate with the whale's overall growth.
An adult whale usually measures in length and in weight. The fluke of this species measures long and the 230 to 360 baleen plates are thought to grow to long, which is longer than that of any other whale by more than a meter. This species is sexually dimorphic as females usually reach lengths of , while males average . There are, however, some specimens that exceed these sizes. In one instance, a female killed off the waters of Pond Inlet in the 1800s allegedly measured . Some estimates put the total maximum length higher at about .
Analysis of hundreds of DNA samples from living whales and from baleen used in vessels, toys, and housing material has shown that Arctic bowhead whales have lost a significant portion of their genetic diversity in the past 500 years. Bowheads originally crossed ice-covered inlets and straits to exchange genes between Atlantic and Pacific populations. This conclusion was derived from analyzing maternal lineage using mitochondrial DNA. Whaling and climatic cooling during the Little Ice Age, from the 16th century to the 19th, is supposed to have reduced the whales' summer habitats, which explains the loss of genetic diversity.
A 2013 discovery has clarified the function of the bowhead's large palatal retial organ. The bulbous ridge of highly vascularized tissue, the corpus cavernosum maxillaris, extends along the centre of the hard plate, forming two large lobes at the rostral palate. The tissue is histologically similar to that of the corpus cavernosum of the mammalian penis. This organ is thought to provide a mechanism of cooling for the whale (which is normally protected from the cold Arctic waters by or more of fat). During physical exertion, the whale must cool itself to prevent hyperthermia (and ultimately brain damage). This organ becomes engorged with blood, and as the whale opens its mouth cold seawater flows over the organ, thus cooling the blood.
In one study, the brain size of two males that measured in total length, were recorded at , respectively. With a gyrencephalic index of 2.32, the brains of the two males were found to exhibit extreme gyrification. Compared to other cetaceans, their brain had a lower level of gyrification in the cerebral cortex, more vertically-aligned gyri, and a relatively dull temporal pole region.
Behaviour
Swimming
The bowhead whale is not a social animal, typically travelling alone or in small pods of up to six. It is able to dive and remain submerged under water for up to an hour. The time spent under water in a single dive is usually limited to 9–18 minutes. The bowhead is not thought to be a deep diver, but can reach a depth down to . It is a slow swimmer, normally travelling around . When fleeing from danger, it can travel at a speed of . During periods of feeding, the average swim speed is increased to .
Feeding
The head of the bowhead whale comprises a large portion of its body length, creating an enormous feeding apparatus. The bowhead whale is a filter feeder, and feeds by swimming forward with its mouth wide open. It has hundreds of overlapping baleen plates consisting of keratin hanging from each side of the upper jaw. The mouth has a large, upturning lip on the lower jaw that helps to reinforce and hold the baleen plates within the mouth. This also prevents buckling or breakage of the plates from the pressure of the water passing through them as the whale advances. To feed, water is filtered through the fine hairs of keratin of the baleen plates, trapping the prey inside near the tongue where it is then swallowed. The diet consists of mostly zooplankton, which includes krill, copepods, mysids, amphipods, and many other crustaceans. About of food are consumed each day. While foraging, bowheads are solitary or occur in groups of two to 10 or more.
Vocalization
Bowhead whales are highly vocal and use low frequency (<1000 Hz) sounds to communicate while travelling, feeding, and socialising. Intense calls for communication and navigation are produced especially during migration season. During breeding season, bowheads make long, complex, variable songs for mating calls. Many tens of distinct songs are sung by a population in a single season. From 2010 through to 2014, near Greenland, 184 distinct songs were recorded from a population of around 300 animals.
Reproduction
Sexual activity occurs between pairs and in boisterous groups of several males and one or two females. Breeding season is observed from March through August; conception is believed to occur primarily in March when song activity is at its highest. Reproduction can begin when a whale is 10 to 15 years old. The gestation period is 13–14 months with females producing a calf once every three to four years. Lactation typically lasts about a year. To survive in the cold water immediately after birth, calves are born with a thick layer of blubber. Within 30 minutes of birth, bowhead calves are able to swim on their own. A newborn calf is typically long, weighs roughly , and grows to within the first year.
Health
Lifespan
Bowhead whales are considered to be the longest-living mammals, living for over 200 years. In May 2007, a specimen caught off the Alaskan coast was discovered with the head of an explosive bomb lance of a model manufactured between 1879 and 1885 lodged in its body, so the whale was probably bomb lanced sometime between those years, and its age at the time of death was estimated at between 115 and 130 years. Spurred by this discovery, scientists measured the ages of other bowhead whales; one specimen was estimated to be 211 years old. Other bowhead whales were estimated to be between 135 and 172 years old. This discovery showed the longevity of the bowhead whale is much greater than originally thought. Researchers at CSIRO, Australia's national science agency, estimated that bowhead whales' maximum natural lifespan is 268 years based on genetic analysis.
Genetic benefits
A greater number of cells present in an organism was once believed to result in greater chances of mutations that cause age-related diseases and cancer. Although the bowhead whale has thousands of times more cells than other mammals, it has a much higher resistance to cancer and aging. In 2015, scientists from the US and UK were able to successfully map the whale's genome. Through comparative analysis, two alleles that could be responsible for the whale's longevity were identified. These two specific gene mutations linked to the bowhead whale's ability to live longer are the ERCC1 gene and the proliferating cell nuclear antigen (PCNA) gene. ERCC1 is linked to DNA repair and increased cancer resistance. PCNA is also important in DNA repair. These mutations enable bowhead whales to better repair DNA damage, allowing for greater resistance to cancer. The whale's genome may also reveal physiological adaptations such as having low metabolic rates compared to other mammals. Changes in the gene UCP1, a gene involved in thermoregulation, can explain differences in the metabolic rates in cells.
Ecology
Range and habitat
The bowhead whale is the only baleen whale to spend its entire life in the Arctic and subarctic waters. The Alaskan population spends the winter months in the southwestern Bering Sea. The group migrates northward in the spring, following openings in the ice, into the Chukchi and Beaufort seas. The whale's range varies depending on climate changes and on the forming/melting of ice.
Historically, bowhead whales' range may have been broader and more southerly than currently thought. Bowheads were abundant around Labrador, Newfoundland (Strait of Belle Isle) and the northern Gulf of St Lawrence until at least the 16th and 17th centuries. It is unclear whether this was due to the colder climate during these periods. The distribution of Balaena spp. during the Pleistocene were far more southerly as fossils have been excavated from Italy and North Carolina, thus could have overlapped between those of Eubalaena based on those locations.
Population
Generally, five stocks of bowhead whales are recognized: 1) the Western Arctic stock in the Bering, Chukchi, and Beaufort Seas, 2) the Hudson Bay and Foxe Basin stock, 3) the Baffin Bay and Davis Strait stock, 4) the Sea of Okhotsk stock, and 5) the Svalbard-Barents Sea stock. However, recent evidence suggests that the Hudson Bay and Foxe Basin stock and the Baffin Bay and Davis Strait stock should be considered one stock based on genetics and movements of tagged whales.
Western Arctic
The Western Arctic bowhead population, also known as the Bering-Chukchi-Beaufort population, has recovered since the commercial harvest of this stock ceased in the early 1900s. A 2019 study estimated that the Western Arctic population was 12,505; although it was lower than the 2011 value of 16,820, the surveyors believed there was no significant decline in 2011–2019 due to the unusual conditions of whale migration and observation in 2019. The yearly growth rate of the Western Arctic bowhead population was 3.7% from 1978 to 2011. These data suggest that the Western Arctic bowhead stock may be near its precommercial whaling level. Migration patterns of this population are being affected by climate change.
Alaskan Natives continue to hunt small numbers of bowhead whales for subsistence purposes. The Alaska Eskimo Whaling Commission co-manages the bowhead subsistence harvest with the National Oceanic and Atmospheric Administration. The Alaskan villages that participate in the bowhead subsistence harvest include Barrow, Point Hope, Point Lay, Wainwright, Nuiqsut, Kaktovik, Gambell, Savoonga, Kivalina, Wales, and Little Diomede. The annual subsistence harvest of the Western Arctic stock has ranged from 14 to 72, amounting to an estimated 0.1-0.5% of the population.
Baffin Bay and Davis Strait
In March 2008, Canada's Department of Fisheries and Oceans stated the previous estimates in the eastern Arctic had undercounted, with a new estimate of 14,400 animals (range 4,800–43,000). These larger numbers correspond to prewhaling estimates, indicating the population has fully recovered. However, if climate change substantially shrinks sea ice, these whales could be threatened by increased shipping traffic.
The status of other populations is less well known. About 1,200 were off West Greenland in 2006, while the Svalbard population may only number in the tens. However, the numbers have been increasing in recent years.
Hudson Bay and Foxe Basin
The Hudson Bay – Foxe Basin population is distinct from the Baffin BayDavis Strait group. The original population size of this local group is unclear, but possibly about 500 to 600 whales annually summered in the northwestern part of the bay in the 1860s. It is likely that the number of whales that actually inhabit Hudson Bay is much smaller than the total population size of this group, but reports from local indigenous people indicate that this population is increasing over decades. Larger portions of the bay are used for summering, while wintering is on a smaller scale. Some animals winter in Hudson Strait, most notably north of Igloolik Island and north eastern Hudson Bay. Distribution patterns in these regions are affected by the presence of orca, and bowheads can disappear from normal ranges in the presence of atypical numbers of orca. Increased mortality caused by orca attack is a possible outcome of climate change, as reduced ice coverage is expected to result in fewer areas that the bowheads can use for shelter from attack. Whaling grounds in the 19th century stretched from Marble Island to Roes Welcome Sound and to Lyon Inlet and Fisher Strait, and whales still migrate through most of these areas.
Distribution within Hudson Bay is mostly restricted to the northwestern part along with Wager Bay, Repulse Bay, Southampton Island (one of two main known summering areas), Frozen Strait, northern Foxe
Basin, and north of Igloolik in summer. Satellite tracking indicates that some portions of the group within the bay do not venture further south than Whale Cove and areas south of Coats and Mansel Islands. Cow – calf pairs and juveniles up to in length make up the majority of summering aggregation in the northern Foxe Basin, while matured males and noncalving females may use the northwestern part of Hudson Bay. Fewer whales also migrate to the west coast of Hudson Bay and Mansel and Ottawa Islands. Bowhead ranges within Hudson Bay are usually considered not to cover southern parts, but at least some whales migrate to locations further south such as Sanikiluaq and Churchill river mouth.
Congregation within Foxe Basin occurs in a well-defined area of north of Igloolik Island to Fury and Hecla Strait and Kapuiviit and Gifford Fiord, and into Gulf of Boothia and Prince Regent Inlet. Northward migrating along western Foxe Basin to eastern side of the basin also occurs in spring.
Sea of Okhotsk
Not much is known about the endangered Sea of Okhotsk population. To learn more about the population, these mammals have been regularly observed near the Shantar Islands, very close to the shore, such as at Ongachan Bay. Several companies provide whale-watching services, which are mostly land-based. According to Russian scientists, this total population likely does not exceed 400 animals. Scientific research on this population was seldom done before 2009, when researchers studying belugas noticed concentrations of bowheads in the study area. Thus, bowheads in the Sea of Okhotsk were once called "forgotten whales" by researchers. The WWF welcomed the creation of a nature sanctuary in the region
Possibly, vagrants from this population occasionally reach into Asian nations such as off Japan or the Korean Peninsula (although this record might be of a right whale). The first documented report of the species in Japanese waters was of a strayed infant () caught in Osaka Bay on 23 June 1969, and the first live sighting was of a juvenile around Shiretoko Peninsula (the southernmost of ice floe range in the Northern Hemisphere) on 21 to 23 June 2015. Fossils have been excavated on Hokkaido, but it is unclear whether the northern coasts of Japan were once included in seasonal or occasional migration ranges.
Genetic studies suggest Okhotsk population share common ancestry with whales in Bering-Chukchi-Beaufort Seas, and repeated mixings had occurred between whales in the two seas.
Svalbard-Barents Sea
The most endangered but historically largest of all bowhead populations is the Svalbard/Spitsbergen population. Occurring normally in Fram Strait, Barents Sea and Severnaya Zemlya along Kara Sea to Laptev Sea and East Siberian Sea regions, these whales were seen in entire coastal regions in European and Russian Arctic, even reaching to Icelandic and Scandinavian coasts and Jan Mayen in Greenland Sea, and west of Cape Farewell and western Greenland coasts. Also, bowheads in this stock were possibly once abundant in areas adjacent to the White Sea region, where few or no animals currently migrate, such as the Kola and Kanin Peninsula. Today, the number of sightings elsewhere is very small, but with increasing regularities with whales having strong regional connections. Whales have also started approaching townships and inhabited areas such as around Longyearbyen. The waters around the marine mammal sanctuary of Franz Josef Land is possibly functioning as the most important habitat for this population.
It is unclear whether this population is a remnant of the historic Svalbard group, recolonized individuals from other stocks, or if a mixing of these two or more stocks has taken place. In 2015, discoveries of the refuge along eastern Greenland where whaling ships could not reach due to ice floes and largest numbers of whales (80–100 individuals) ever sighted between Spitsbergen and Greenland indicate that more whales than previously considered survived whaling periods, and flows from the other populations are possible.
Possible moulting area off Baffin Island
During expeditions by a tour operator 'Arctic Kingdom', a large group of bowheads seemingly involved in courtship activities was discovered in very shallow bays south of Qikiqtarjuaq in 2012. Floating skins and rubbing behaviours at sea bottom indicated possible moulting had taken place. Moulting behaviours had never or seldom been documented for this species before. This area is an important habitat for whales that were observed to be relatively active and to interact with humans positively, or to rest on sea floors. These whales belong to Davis Strait stock.
Isabella Bay in Niginganiq National Wildlife Area is the first wildlife sanctuary in the world to be designed specially for bowhead whales. However, moultings have not been recorded in this area due to environmental factors.
Predation
In 1978 the International Whaling Commission (IWC) introduced a hunting strike quota for the Bering-Chukchi-Beaufort Sea (BCB) bowhead. The quota has remained at 67 strikes per year since 1998 and represents about 0.5 percent of BCB population. The population of bowheads in West Greenland and Canada is estimated to be 6,000 and rising, and hunts in this are minimal (<0.001 percent). Both stocks are rising, and the indigenous hunts seem to be self-sustaining. Orca are also known predators. There is no consensus on the number of deaths by orca. Bowheads seek the ice and shallow waters' safety when threatened by orca. The Inuit have a traditional word for this behavior to give historical context that this is not a new phenomenon. Global warming is increasing the frequency that orca are observed in the far north. A once-rare event, orca are now seen more frequently.
There are no reports of attacks on bowheads by sharks.
Whaling
The bowhead whale has been hunted for blubber, meat, oil, bones, and baleen. Like the right whale, it swims slowly, and floats after death, making it ideal for whaling. Before commercial whaling, they were estimated to number 50,000. Paleo-Eskimo sites indicate bowhead whales were eaten in sites from perhaps 4000 BC. Inuit people near the Pacific developed specific hunting tools, with the whales providing food and fuel.
Commercial bowhead whaling began in the 16th century when the Basques killed them as they migrated south through the Strait of Belle Isle in the fall and early winter. In 1611, the first whaling expedition sailed to Spitsbergen. The whaling settlement Smeerenburg was founded on Spitsbergen in 1619. By midcentury, the population(s) there had practically been wiped out, forcing whalers to voyage into the "West Ice"—the pack ice off Greenland's east coast. By 1719, they had reached the Davis Strait, and by the first quarter of the 19th century, Baffin Bay.
In the North Pacific, the first bowheads were taken off the eastern coast of Kamchatka by the Danish whaleship Neptun, Captain Thomas Sodring, in 1845. In 1847, the first bowheads were caught in the Sea of Okhotsk, and the following year, Captain Thomas Welcome Roys, in the bark Superior, of Sag Harbor, caught the first bowheads in the Bering Strait region. By 1849, 50 ships were hunting bowheads in each area; in the Bering Strait, 500 whales were killed that year, and that number jumped to more than 2000 in 1850. By 1852, 220 ships were cruising around the Bering Strait region, which killed over 2,600 whales. Between 1854 and 1857, the fleet shifted to the Sea of Okhotsk, where 100–160 ships cruised annually. During 1858–1860, the ships shifted back to the Bering Strait region, where the majority of the fleet cruised during the summer until the early 20th century. An estimated 18,600 bowheads were killed in the Bering Strait region between 1848 and 1914, with 60% of the total being reached within the first two decades. An estimated 18,000 bowheads were killed in the Sea of Okhotsk during 1847–1867, 80% in the first decade.
Bowheads were first taken along the pack ice in the northeastern Sea of Okhotsk, then in Tausk Bay and Northeast Gulf (Shelikhov Gulf). Soon, ships expanded to the west, catching them around Iony Island and then around the Shantar Islands. In the Western Arctic, they mainly caught them in the Anadyr Gulf, the Bering Strait, and around St. Lawrence Island. They later spread to the western Beaufort Sea (1854) and the Mackenzie River delta (1889).
Commercial whaling, the principal cause of the population decline, is over. Bowhead whales are now hunted on a subsistence level by native peoples of North America.
In 2024, the Inuit hunters of Aklavik, Northwest Territories were permitted to hunt and kill one bowhead whale to distribute the whale meat, an important part of Inuit cuisine, to Inuvialuit and Gwich'in communities in the region.
Conservation
The bowhead is listed in Appendix I by CITES. While the global population is thought to be secure, thus assigned "least concern" status, some populations are listed by the National Marine Fisheries Service as "endangered" under the auspices of the United States' Endangered Species Act. The IUCN Red List data are:
Svalbard population – critically endangered
Sea of Okhotsk subpopulation – endangered
Baffin Bay-Davis Strait stock – endangered
Hudson Bay-Foxe Basin stock – vulnerable (estimated to be 1,026 individuals in 2005 by DFO)
Bering-Chukchi-Beaufort stock – lower risk – conservation dependent
The Alaska Department of Fish and Game and the USA government list the bowhead whale as federally endangered.
The bowhead whale is listed in Appendix I of the Convention on the Conservation of Migratory Species of Wild Animals (CMS), as this species has been categorized as being in danger of extinction throughout all or a significant proportion of its range. CMS Parties strive towards strictly protecting these animals, conserving or restoring the places where they live, mitigating obstacles to migration, and controlling other factors that might endanger them.
| Biology and health sciences | Baleen whales | Animals |
20739459 | https://en.wikipedia.org/wiki/Dissolved%20load | Dissolved load | Dissolved load is the portion of a stream's total sediment load that is carried in solution, especially ions from chemical weathering. It is a major contributor to the total amount of material removed from a river's drainage basin, along with suspended load and bed load. The amount of material carried as dissolved load is typically much smaller than the suspended load, though this is not always the case, particularly when the available river flow is mostly harnessed for purposes such as irrigation or industrial uses. Dissolved load comprises a significant portion of the total material flux out of a landscape, and its composition is important in regulating the chemistry and biology of the stream water.
The dissolved load is primarily controlled by the rate of chemical weathering, which depends on climate and weather conditions such as moisture and temperature. Dissolved load has many useful applications within the field of geology, including erosion, denudation, and reconstructing climate in the past.
Measurement techniques
Dissolved load is typically measured by taking samples of water from a river and running various scientific tests on them. First, the pH, conductivity, and bicarbonate alkalinity of the sample are measured. Next, samples are filtered to remove any suspended sediments and preserved with chloroform to prevent growth of microorganisms, while the others are acidified with hydrochloric acid added to keep dissolved ions from precipitating out of solution. Then, various chemical tests are applied to determine the concentration of each solute. For example, the concentrations of sodium and potassium ions can be determined by flame photometry, while the calcium and magnesium ion concentrations can be determined by atomic absorption spectrophotometry.
Applications
Reconstructing climate
Dissolved load can provide valuable information about the rate of soil formation and other processes of chemical erosion. In particular, the mass balance between the dissolved load and solid phase is helpful in determining surface dynamics. In addition, dissolved load can be used to reconstruct the climate of the Earth in the past. This is because chemical weathering is the major contributor to the dissolved load of a stream. The chemical weathering of silicate rocks is the primary sink for carbon dioxide in the atmosphere, because atmospheric carbon dioxide is converted into carbonate rocks in the carbonate–silicate cycle. Carbon dioxide concentrations are the primary control of the greenhouse effect, which determines the temperature of the Earth.
Denudation
Denudation is the process of wearing away the top layers of Earth's landscape. Because the denudation rate is normally too low to measure directly, it can be indirectly determined by measuring the sediment load of the streams that drain the area in question. This is possible because any material that passes through a certain point on a stream is guaranteed to have come from somewhere in the stream's drainage basin upstream of that point. As topographic relief increases, the dissolved load's contribution to the total stream load decreases due to the fact that on steeper surfaces, rain is less likely to infiltrate the rocks, leading to less chemical weathering, which decreases the dissolved load.
Salt export
The process of carrying salts by water to the sea or a land-locked lake from a river basin is called salt export. When adequate salt export is not occurring, the river basin area gradually converts into saline soils and/or alkali soils, particularly in lower reaches.
Dissolved loads of selected major rivers
| Physical sciences | Sedimentology | Earth science |
20741014 | https://en.wikipedia.org/wiki/Pentagon | Pentagon | In geometry, a pentagon () is any five-sided polygon or 5-gon. The sum of the internal angles in a simple pentagon is 540°.
A pentagon may be simple or self-intersecting. A self-intersecting regular pentagon (or star pentagon) is called a pentagram.
Regular pentagons
A regular pentagon has Schläfli symbol {5} and interior angles of 108°.
A regular pentagon has five lines of reflectional symmetry, and rotational symmetry of order 5 (through 72°, 144°, 216° and 288°). The diagonals of a convex regular pentagon are in the golden ratio to its sides. Given its side length its height (distance from one side to the opposite vertex), width (distance between two farthest separated points, which equals the diagonal length ) and circumradius are given by:
The area of a convex regular pentagon with side length is given by
If the circumradius of a regular pentagon is given, its edge length is found by the expression
and its area is
since the area of the circumscribed circle is the regular pentagon fills approximately 0.7568 of its circumscribed circle.
Derivation of the area formula
The area of any regular polygon is:
where P is the perimeter of the polygon, and r is the inradius (equivalently the apothem). Substituting the regular pentagon's values for P and r gives the formula
with side length t.
Inradius
Similar to every regular convex polygon, the regular convex pentagon has an inscribed circle. The apothem, which is the radius r of the inscribed circle, of a regular pentagon is related to the side length t by
Chords from the circumscribed circle to the vertices
Like every regular convex polygon, the regular convex pentagon has a circumscribed circle. For a regular pentagon with successive vertices A, B, C, D, E, if P is any point on the circumcircle between points B and C, then PA + PD = PB + PC + PE.
Point in plane
For an arbitrary point in the plane of a regular pentagon with circumradius , whose distances to the centroid of the regular pentagon and its five vertices are and
respectively, we have
If are the distances from the vertices of a regular pentagon to any point on its circumcircle, then
Geometrical constructions
The regular pentagon is constructible with compass and straightedge, as 5 is a Fermat prime. A variety of methods are known for constructing a regular pentagon. Some are discussed below.
Richmond's method
One method to construct a regular pentagon in a given circle is described by Richmond and further discussed in Cromwell's Polyhedra.
The top panel shows the construction used in Richmond's method to create the side of the inscribed pentagon. The circle defining the pentagon has unit radius. Its center is located at point C and a midpoint M is marked halfway along its radius. This point is joined to the periphery vertically above the center at point D. Angle CMD is bisected, and the bisector intersects the vertical axis at point Q. A horizontal line through Q intersects the circle at point P, and chord PD is the required side of the inscribed pentagon.
To determine the length of this side, the two right triangles DCM and QCM are depicted below the circle. Using Pythagoras' theorem and two sides, the hypotenuse of the larger triangle is found as . Side h of the smaller triangle then is found using the half-angle formula:
where cosine and sine of ϕ are known from the larger triangle. The result is:
If DP is truly the side of a regular pentagon, , so DP = 2 cos(54°), QD = DP cos(54°) = 2cos2(54°), and CQ = 1 − 2cos2(54°), which equals −cos(108°) by the cosine double angle formula. This is the cosine of 72°, which equals as desired.
Carlyle circles
The Carlyle circle was invented as a geometric method to find the roots of a quadratic equation. This methodology leads to a procedure for constructing a regular pentagon. The steps are as follows:
Draw a circle in which to inscribe the pentagon and mark the center point O.
Draw a horizontal line through the center of the circle. Mark the left intersection with the circle as point B.
Construct a vertical line through the center. Mark one intersection with the circle as point A.
Construct the point M as the midpoint of O and B.
Draw a circle centered at M through the point A. Mark its intersection with the horizontal line (inside the original circle) as the point W and its intersection outside the circle as the point V.
Draw a circle of radius OA and center W. It intersects the original circle at two of the vertices of the pentagon.
Draw a circle of radius OA and center V. It intersects the original circle at two of the vertices of the pentagon.
The fifth vertex is the rightmost intersection of the horizontal line with the original circle.
Steps 6–8 are equivalent to the following version, shown in the animation:
6a. Construct point F as the midpoint of O and W.
7a. Construct a vertical line through F. It intersects the original circle at two of the vertices of the pentagon. The third vertex is the rightmost intersection of the horizontal line with the original circle.
8a. Construct the other two vertices using the compass and the length of the vertex found in step 7a.
Euclid's method
A regular pentagon is constructible using a compass and straightedge, either by inscribing one in a given circle or constructing one on a given edge. This process was described by Euclid in his Elements circa 300 BC.
Physical construction methods
A regular pentagon may be created from just a strip of paper by tying an overhand knot into the strip and carefully flattening the knot by pulling the ends of the paper strip. Folding one of the ends back over the pentagon will reveal a pentagram when backlit.
Construct a regular hexagon on stiff paper or card. Crease along the three diameters between opposite vertices. Cut from one vertex to the center to make an equilateral triangular flap. Fix this flap underneath its neighbor to make a pentagonal pyramid. The base of the pyramid is a regular pentagon.
Symmetry
The regular pentagon has Dih5 symmetry, order 10. Since 5 is a prime number there is one subgroup with dihedral symmetry: Dih1, and 2 cyclic group symmetries: Z5, and Z1.
These 4 symmetries can be seen in 4 distinct symmetries on the pentagon. John Conway labels these by a letter and group order. Full symmetry of the regular form is r10 and no symmetry is labeled a1. The dihedral symmetries are divided depending on whether they pass through vertices (d for diagonal) or edges (p for perpendiculars), and i when reflection lines path through both edges and vertices. Cyclic symmetries in the middle column are labeled as g for their central gyration orders.
Each subgroup symmetry allows one or more degrees of freedom for irregular forms. Only the g5 subgroup has no degrees of freedom but can be seen as directed edges.
Regular pentagram
A pentagram or pentangle is a regular star pentagon. Its Schläfli symbol is {5/2}. Its sides form the diagonals of a regular convex pentagon – in this arrangement the sides of the two pentagons are in the golden ratio.
Equilateral pentagons
An equilateral pentagon is a polygon with five sides of equal length. However, its five internal angles can take a range of sets of values, thus permitting it to form a family of pentagons. In contrast, the regular pentagon is unique up to similarity, because it is equilateral and it is equiangular (its five angles are equal).
Cyclic pentagons
A cyclic pentagon is one for which a circle called the circumcircle goes through all five vertices. The regular pentagon is an example of a cyclic pentagon. The area of a cyclic pentagon, whether regular or not, can be expressed as one fourth the square root of one of the roots of a septic equation whose coefficients are functions of the sides of the pentagon.
There exist cyclic pentagons with rational sides and rational area; these are called Robbins pentagons. It has been proven that the diagonals of a Robbins pentagon must be either all rational or all irrational, and it is conjectured that all the diagonals must be rational.
General convex pentagons
For all convex pentagons with sides and diagonals , the following inequality holds:
.
Pentagons in tiling
A regular pentagon cannot appear in any tiling of regular polygons. First, to prove a pentagon cannot form a regular tiling (one in which all faces are congruent, thus requiring that all the polygons be pentagons), observe that (where 108° Is the interior angle), which is not a whole number; hence there exists no integer number of pentagons sharing a single vertex and leaving no gaps between them. More difficult is proving a pentagon cannot be in any edge-to-edge tiling made by regular polygons:
The maximum known packing density of a regular pentagon is , achieved by the double lattice packing shown. In a preprint released in 2016, Thomas Hales and Wöden Kusner announced a proof that this double lattice packing of the regular pentagon (known as the "pentagonal ice-ray" Chinese lattice design, dating from around 1900) has the optimal density among all packings of regular pentagons in the plane.
There are no combinations of regular polygons with 4 or more meeting at a vertex that contain a pentagon. For combinations with 3, if 3 polygons meet at a vertex and one has an odd number of sides, the other 2 must be congruent. The reason for this is that the polygons that touch the edges of the pentagon must alternate around the pentagon, which is impossible because of the pentagon's odd number of sides. For the pentagon, this results in a polygon whose angles are all . To find the number of sides this polygon has, the result is , which is not a whole number. Therefore, a pentagon cannot appear in any tiling made by regular polygons.
There are 15 classes of pentagons that can monohedrally tile the plane. None of the pentagons have any symmetry in general, although some have special cases with mirror symmetry.
Pentagons in polyhedra
Pentagons in nature
Plants
Animals
Minerals
Other examples
| Mathematics | Two-dimensional space | null |
12773567 | https://en.wikipedia.org/wiki/Arboreal%20locomotion | Arboreal locomotion | Arboreal locomotion is the locomotion of animals in trees. In habitats in which trees are present, animals have evolved to move in them. Some animals may scale trees only occasionally (scansorial), but others are exclusively arboreal. The habitats pose numerous mechanical challenges to animals moving through them and lead to a variety of anatomical, behavioral and ecological consequences as well as variations throughout different species. Furthermore, many of these same principles may be applied to climbing without trees, such as on rock piles or mountains.
Some animals are exclusively arboreal in habitat, such as tree snails.
Biomechanics
Arboreal habitats pose numerous mechanical challenges to animals moving in them, which have been solved in diverse ways. These challenges include moving on narrow branches, moving up and down inclines, balancing, crossing gaps, and dealing with obstructions.
Diameter
Moving along narrow surfaces, such as a branch of a tree, can create special difficulties for animals who are not adapted to deal with balancing on small diameter substrates. During locomotion on the ground, the location of the center of mass may swing from side to side. But during arboreal locomotion, this would result in the center of mass moving beyond the edge of the branch, resulting in a tendency to topple over and fall. Not only do some arboreal animals have to be able to move on branches of varying diameter, but they also have to eat on these branches, resulting in the need for the ability to balance while using their hands to feed themselves. This resulted in various types of grasping such as pedal grasping in order to clamp themselves onto small branches for better balance.
Incline
Branches are frequently oriented at an angle to gravity in arboreal habitats, including being vertical, which poses special problems. As an animal moves up an inclined branch, it must fight the force of gravity to raise its body, making the movement more difficult. To get past this difficulty, many animals have to grasp the substrate with all four limbs and increase the frequency of their gait sequence. Conversely, as the animal descends, it must also fight gravity to control its descent and prevent falling. Descent can be particularly problematic for many animals, and highly arboreal species often have specialized methods for controlling their descent. One way animals prevent falling while descending is to increase the amount of contact their limbs are making with the substrate to increase friction and braking power.
Balance
Due to the height of many branches and the potentially disastrous consequences of a fall, balance is of primary importance to arboreal animals. On horizontal and gently sloped branches, the primary problem is tipping to the side due to the narrow base of support. The narrower the branch, the greater the difficulty in balancing a given animal faces. On steep and vertical branches, tipping becomes less of an issue, and pitching backwards or slipping downwards becomes the most likely failure. In this case, large-diameter branches pose a greater challenge since the animal cannot place its forelimbs closer to the center of the branch than its hindlimbs.
Crossing gaps
Some arboreal animals need to be able to move from tree to tree in order to find food and shelter. To be able to get from tree to tree, animals have evolved various adaptations. In some areas trees are close together and can be crossed by simple brachiation. In other areas, trees are not close together and animals need to have specific adaptations to jump far distances or glide.
Obstructions
Arboreal habitats often contain many obstructions, both in the form of branches emerging from the one being moved on and other branches impinging on the space the animal needs to move through. These obstructions may impede locomotion, or may be used as additional contact points to enhance it. While obstructions tend to impede limbed animals, they benefit snakes by providing anchor points.
Anatomical specializations
Arboreal organisms display many specializations for dealing with the mechanical challenges of moving through their habitats.
Arboreal animals frequently have elongated limbs that help them cross gaps, reach fruit or other resources, test the firmness of support ahead, and in some cases, to brachiate. However, some species of lizard have reduced limb size that helps them avoid limb movement being obstructed by impinging branches.
Many arboreal species, such as howler monkeys, green tree pythons, emerald tree boas, chameleons, silky anteaters, spider monkeys, and possums, use prehensile tails to grasp branches. In the spider monkey and crested gecko, the tip of the tail has either a bare patch or adhesive pad, which provides increased friction.
Claws can be used to interact with rough substrates and re-orient the direction of the force the animal applies. This is what allows squirrels to climb tree trunks that are so large as to be essentially flat, from the perspective of such a small animal. However, claws can interfere with an animal's ability to grasp very small branches, as they may wrap too far around and prick the animal's own paw.
Adhesion is an alternative to claws, which works best on smooth surfaces. Wet adhesion is common in tree frogs and arboreal salamanders, and functions either by suction or by capillary adhesion. Dry adhesion is best typified by the specialized toes of geckos, which use van der Waals forces to adhere to many substrates, even glass.
Frictional gripping is used by primates, relying upon hairless fingertips. Squeezing the branch between the fingertips generates a frictional force that holds the animal's hand to the branch. However, this type of grip depends upon the angle of the frictional force; thus upon the diameter of the branch, with larger branches resulting in reduced gripping ability. Animals other than primates that use gripping in climbing include the chameleon, which has mitten-like grasping feet, and many birds that grip branches in perching or moving about.
To control descent, especially down large diameter branches, some arboreal animals such as squirrels have evolved highly mobile ankle joints that permit rotating the foot into a 'reversed' posture. This allows the claws to hook into the rough surface of the bark, opposing the force of gravity.
Many arboreal species lower their center of mass to reduce pitching and toppling movement when climbing. This may be accomplished by postural changes, altered body proportions, or smaller size.
Small size provides many advantages to arboreal species: such as increasing the relative size of branches to the animal, lower center of mass, increased stability, lower mass (allowing movement on smaller branches), and the ability to move through more cluttered habitat. Size relating to weight affects gliding animals such as the reduced weight per snout-vent length for 'flying' frogs.
Some species of primate, bat, and all species of sloth achieve passive stability by hanging beneath the branch. Both pitching and tipping become irrelevant, as the only method of failure would be losing their grip.
Behavioral specializations
Arboreal species have behaviors specialized for moving in their habitats, most prominently in terms of posture and gait. Specifically, arboreal mammals take longer steps, extend their limbs further forwards and backwards during a step, adopt a more 'crouched' posture to lower their center of mass, and use a diagonal sequence gait.
Brachiation
Brachiation is a specialized form of arboreal locomotion, used by primates to move very rapidly while hanging beneath branches. Arguably the epitome of arboreal locomotion, it involves swinging with the arms from one handhold to another. Only a few species are brachiators, and all of these are primates; it is a major means of locomotion among spider monkeys and gibbons, and is occasionally used by female orangutans. Gibbons are the experts of this mode of locomotion, swinging from branch to branch distances of up to 15 m (50 ft), and traveling at speeds of as much as .
Gliding and parachuting
To bridge gaps between trees, many animals such as the flying squirrel have adapted membranes, such as patagia for gliding flight. Some animals can slow their descent in the air using a method known as parachuting, such as Rhacophorus (a "flying frog" species) that has adapted toe membranes allowing it to fall more slowly after leaping from trees.
Limbless climbing
Many species of snake are highly arboreal, and some have evolved specialized musculature for this habitat. While moving in arboreal habitats, snakes move slowly along bare branches using a specialized form of concertina locomotion, but when secondary branches emerge from the branch being moved on, snakes use lateral undulation, a much faster mode. As a result, snakes perform best on small perches in cluttered environments, while limbed organisms seem to do best on large perches in uncluttered environments.
Evolutionary history
The earliest known climbing tetrapod is the varanopid amniote Eoscansor from the Late Carboniferous (Pennsylvanian) of North America which is clearly specialised with adaptations for grasping, likely onto tree trunks. Suminia, a anomodont synapsid from Russia dating to the Late Permian, about 260 million years ago, was also likely a specialised climber.
| Biology and health sciences | Ethology | Biology |
10490785 | https://en.wikipedia.org/wiki/Evolution%20of%20the%20horse | Evolution of the horse | The evolution of the horse, a mammal of the family Equidae, occurred over a geologic time scale of 50 million years, transforming the small, dog-sized, forest-dwelling Eohippus into the modern horse. Paleozoologists have been able to piece together a more complete outline of the evolutionary lineage of the modern horse than of any other animal. Much of this evolution took place in North America, where horses originated but became extinct about 10,000 years ago, before being reintroduced in the 15th century.
The horse belongs to the order Perissodactyla (odd-toed ungulates), the members of which all share hooved feet and an odd number of toes on each foot, as well as mobile upper lips and a similar tooth structure. This means that horses share a common ancestry with tapirs and rhinoceroses. The perissodactyls arose in the late Paleocene, less than 10 million years after the Cretaceous–Paleogene extinction event. This group of animals appears to have been originally specialized for life in tropical forests, but whereas tapirs and, to some extent, rhinoceroses, retained their jungle specializations, modern horses are adapted to life in the climatic conditions of the steppes, which are drier and much harsher than forests or jungles. Other species of Equus are adapted to a variety of intermediate conditions.
The early ancestors of the modern horse walked on several spread-out toes, an accommodation to life spent walking on the soft, moist ground of primeval forests. As grass species began to appear and flourish, the equids' diets shifted from foliage to silicate-rich grasses; the increased wear on teeth selected for increases in the size and durability of teeth. At the same time, as the steppes began to appear, selection favored increase in speed to outrun predators. This ability was attained by lengthening of limbs and the lifting of some toes from the ground in such a way that the weight of the body was gradually placed on one of the longest toes, the third.
History of research
Wild horses have been known since prehistory from central Asia to Europe, with domestic horses and other equids being distributed more widely in the Old World, but no horses or equids of any type were found in the New World when European explorers reached the Americas. When the Spanish colonists brought domestic horses from Europe, beginning in 1493, escaped horses quickly established large feral herds. In the 1760s, the early naturalist Buffon suggested this was an indication of inferiority of the New World fauna, but later reconsidered this idea. William Clark's 1807 expedition to Big Bone Lick found "leg and foot bones of the Horses", which were included with other fossils sent to Thomas Jefferson and evaluated by the anatomist Caspar Wistar, but neither commented on the significance of this find.
The first Old World equid fossil was found in the gypsum quarries in Montmartre, Paris, in the 1820s. The tooth was sent to the Paris Conservatory, where it was identified by Georges Cuvier, who identified it as a browsing equine related to the tapir. His sketch of the entire animal matched later skeletons found at the site.
During the Beagle survey expedition, the young naturalist Charles Darwin had remarkable success with fossil hunting in Patagonia. On 10 October 1833, at Santa Fe, Argentina, he was "filled with astonishment" when he found a horse's tooth in the same stratum as fossil giant armadillos, and wondered if it might have been washed down from a later layer, but concluded this was "not very probable". After the expedition returned in 1836, the anatomist Richard Owen confirmed the tooth was from an extinct species, which he subsequently named Equus curvidens, and remarked, "This evidence of the former existence of a genus, which, as regards South America, had become extinct, and has a second time been introduced into that Continent, is not one of the least interesting fruits of Mr. Darwin's palæontological discoveries."
In 1848, a study On the fossil horses of America by Joseph Leidy systematically examined Pleistocene horse fossils from various collections, including that of the Academy of Natural Sciences, and concluded at least two ancient horse species had existed in North America: Equus curvidens and another, which he named Equus americanus. A decade later, however, he found the latter name had already been taken and renamed it Equus complicatus. In the same year, he visited Europe and was introduced by Owen to Darwin.
The original sequence of species believed to have evolved into the horse was based on fossils discovered in North America in 1879 by paleontologist Othniel Charles Marsh. The sequence, from Eohippus to the modern horse (Equus), was popularized by Thomas Huxley and became one of the most widely known examples of a clear evolutionary progression. The horse's evolutionary lineage became a common feature of biology textbooks, and the sequence of transitional fossils was assembled by the American Museum of Natural History into an exhibit that emphasized the gradual, "straight-line" evolution of the horse.
Since then, as the number of equid fossils has increased, the actual evolutionary progression from Eohippus to Equus has been discovered to be much more complex and multibranched than was initially supposed. The straight, direct progression from the former to the latter has been replaced by a more elaborate model with numerous branches in different directions, of which the modern horse is only one of many. George Gaylord Simpson in 1951 first recognized that the modern horse was not the "goal" of the entire lineage of equids, but is simply the only genus of the many horse lineages to survive.
Detailed fossil information on the distribution and rate of change of new equid species has also revealed that the progression between species was not as smooth and consistent as was once believed. Although some transitions, such as that of Dinohippus to Equus, were indeed gradual progressions, a number of others, such as that of Epihippus to Mesohippus, were relatively abrupt in geologic time, taking place over only a few million years. Both anagenesis (gradual change in an entire population's gene frequency) and cladogenesis (a population "splitting" into two distinct evolutionary branches) occurred, and many species coexisted with "ancestor" species at various times. The change in equids' traits was also not always a "straight line" from Eohippus to Equus: some traits reversed themselves at various points in the evolution of new equid species, such as size and the presence of facial fossae, and only in retrospect can certain evolutionary trends be recognized.
Before odd-toed ungulates
Phenacodontidae
Phenacodontidae is the most recent family in the order Condylarthra believed to be ancestral to the odd-toed ungulates. It contains the genera Almogaver, Copecion, Ectocion, Eodesmatodon, Meniscotherium, Ordathspidotherium, Phenacodus and Pleuraspidotherium. The family lived from the Early Paleocene to the Middle Eocene in Europe and were about the size of a sheep, with tails making slightly less than half of the length of their bodies and unlike their ancestors, good running skills.
Eocene and Oligocene: early equids
Eohippus
Eohippus appeared in the Ypresian (early Eocene), about 52 mya (million years ago). It was an animal approximately the size of a fox (250–450 mm in height), with a relatively short head and neck and a springy, arched back. It had 44 low-crowned teeth, in the typical arrangement of an omnivorous, browsing mammal: three incisors, one canine, four premolars, and three molars on each side of the jaw. Its molars were uneven, dull, and bumpy, and used primarily for grinding foliage. The cusps of the molars were slightly connected in low crests. Eohippus browsed on soft foliage and fruit, probably scampering between thickets in the mode of a modern muntjac. It had a small brain, and possessed especially small frontal lobes.
Its limbs were long relative to its body, already showing the beginnings of adaptations for running. However, all of the major leg bones were unfused, leaving the legs flexible and rotatable. Its wrist and hock joints were low to the ground. The forelimbs had developed five toes, of which four were equipped with small proto-hooves; the large fifth "toe-thumb" was off the ground. The hind limbs had small hooves on three out of the five toes, whereas the vestigial first and fifth toes did not touch the ground. Its feet were padded, much like a dog's, but with the small hooves in place of claws.
For a span of about 20 million years, Eohippus thrived with few significant evolutionary changes. The most significant change was in the teeth, which began to adapt to its changing diet, as these early Equidae shifted from a mixed diet of fruits and foliage to one focused increasingly on browsing foods. During the Eocene, an Eohippus species (most likely Eohippus angustidens) branched out into various new types of Equidae. Thousands of complete, fossilized skeletons of these animals have been found in the Eocene layers of North American strata, mainly in the Wind River basin in Wyoming. Similar fossils have also been discovered in Europe, such as Propalaeotherium (which is not considered ancestral to the modern horse).
Orohippus
Approximately 50 million years ago, in the early-to-middle Eocene, Eohippus smoothly transitioned into Orohippus through a gradual series of changes. Although its name means "mountain horse", Orohippus was not a true horse and did not live in the mountains. It resembled Eohippus in size, but had a slimmer body, an elongated head, slimmer forelimbs, and longer hind legs, all of which are characteristics of a good jumper. Although Orohippus was still pad-footed, the vestigial outer toes of Eohippus were not present in Orohippus; there were four toes on each fore leg, and three on each hind leg.
The most dramatic change between Eohippus and Orohippus was in the teeth: the first of the premolar teeth was dwarfed, the last premolar shifted in shape and function into a molar, and the crests on the teeth became more pronounced. Both of these factors increased the grinding ability of the teeth of Orohippus; the change suggest selection imposed by increased toughness of Orohippus plant diet.
Epihippus
In the mid-Eocene, about 47 million years ago, Epihippus, a genus which continued the evolutionary trend of increasingly efficient grinding teeth, evolved from Orohippus. Epihippus had five grinding, low-crowned cheek teeth with well-formed crests. A late species of Epihippus, sometimes referred to as Duchesnehippus intermedius, had teeth similar to Oligocene equids, although slightly less developed. Whether Duchesnehippus was a subgenus of Epihippus or a distinct genus is disputed. Epihippus was only 2 feet tall.
Mesohippus
In the late Eocene and the early stages of the Oligocene epoch (32–24 mya), the climate of North America became drier, and the earliest grasses began to evolve. The forests were yielding to flatlands, home to grasses and various kinds of brush. In a few areas, these plains were covered in sand, creating the type of environment resembling the present-day prairies.
In response to the changing environment, the then-living species of Equidae also began to change. In the late Eocene, they began developing tougher teeth and becoming slightly larger and leggier, allowing for faster running speeds in open areas, and thus for evading predators in nonwooded areas. About 40 mya, Mesohippus ("middle horse") suddenly developed in response to strong new selective pressures to adapt, beginning with the species Mesohippus celer and soon followed by Mesohippus westoni.
In the early Oligocene, Mesohippus was one of the more widespread mammals in North America. It walked on three toes on each of its front and hind feet (the first and fifth toes remained, but were small and not used in walking). The third toe was stronger than the outer ones, and thus more weighted; the fourth front toe was diminished to a vestigial nub. Judging by its longer and slimmer limbs, Mesohippus was an agile animal.
Mesohippus was slightly larger than Epihippus, about 610 mm (24 in) at the shoulder. Its back was less arched, and its face, snout, and neck were somewhat longer. It had significantly larger cerebral hemispheres, and had a small, shallow depression on its skull called a fossa, which in modern horses is quite detailed. The fossa serves as a useful marker for identifying an equine fossil's species. Mesohippus had six grinding "cheek teeth", with a single premolar in front—a trait all descendant Equidae would retain. Mesohippus also had the sharp tooth crests of Epihippus, improving its ability to grind down tough vegetation.
Miohippus
Around 36 million years ago, soon after the development of Mesohippus, Miohippus ("lesser horse") emerged, the earliest species being Miohippus assiniboiensis. As with Mesohippus, the appearance of Miohippus was relatively abrupt, though a few transitional fossils linking the two genera have been found. Mesohippus was once believed to have anagenetically evolved into Miohippus by a gradual series of progressions, but new evidence has shown its evolution was cladogenetic: a Miohippus population split off from the main genus Mesohippus, coexisted with Mesohippus for around four million years, and then over time came to replace Mesohippus.
Miohippus was significantly larger than its predecessors, and its ankle joints had subtly changed. Its facial fossa was larger and deeper, and it also began to show a variable extra crest in its upper cheek teeth, a trait that became a characteristic feature of equine teeth.
Miohippus ushered in a major new period of diversification in Equidae.
Miocene and Pliocene: true equines
Kalobatippus
The forest-suited form was Kalobatippus (or Miohippus intermedius, depending on whether it was a new genus or species), whose second and fourth front toes were long, well-suited to travel on the soft forest floors. Kalobatippus probably gave rise to Anchitherium, which travelled to Asia via the Bering Strait land bridge, and from there to Europe. In both North America and Eurasia, larger-bodied genera evolved from Anchitherium: Sinohippus in Eurasia and Hypohippus and Megahippus in North America. Hypohippus became extinct by the late Miocene.
Parahippus
The Miohippus population that remained on the steppes is believed to be ancestral to Parahippus, a North American animal about the size of a small pony, with a prolonged skull and a facial structure resembling the horses of today. Its third toe was stronger and larger, and carried the main weight of the body. Its four premolars resembled the molar teeth; the first were small and almost nonexistent. The incisor teeth, like those of its predecessors, had a crown (like human incisors); however, the top incisors had a trace of a shallow crease marking the beginning of the core/cup.
Merychippus
In the middle of the Miocene epoch, the grazer Merychippus flourished. It had wider molars than its predecessors, which are believed to have been used for crunching the hard grasses of the steppes. The hind legs, which were relatively short, had side toes equipped with small hooves, but they probably only touched the ground when running. Merychippus radiated into at least 19 additional grassland species.
Hipparion
Three lineages within Equidae are believed to be descended from the numerous varieties of Merychippus: Hipparion, Protohippus and Pliohippus. The most different from Merychippus was Hipparion, mainly in the structure of tooth enamel: in comparison with other Equidae, the inside, or tongue side, had a completely isolated parapet. A complete and well-preserved skeleton of the North American Hipparion shows an animal the size of a small pony. They were very slim, rather like antelopes, and were adapted to life on dry prairies. On its slim legs, Hipparion had three toes equipped with small hooves, but the side toes did not touch the ground.
In North America, Hipparion and its relatives (Cormohipparion, Nannippus, Neohipparion, and Pseudhipparion), proliferated into many kinds of equids, at least one of which managed to migrate to Asia and Europe during the Miocene epoch. (European Hipparion differs from American Hipparion in its smaller body size – the best-known discovery of these fossils was near Athens.)
Pliohippus
Pliohippus arose from Callippus in the middle Miocene, around 12 mya. It was very similar in appearance to Equus, though it had two long extra toes on both sides of the hoof, externally barely visible as callused stubs. The long and slim limbs of Pliohippus reveal a quick-footed steppe animal.
Until recently, Pliohippus was believed to be the ancestor of present-day horses because of its many anatomical similarities. However, though Pliohippus was clearly a close relative of Equus, its skull had deep facial fossae, whereas Equus had no fossae at all. Additionally, its teeth were strongly curved, unlike the very straight teeth of modern horses. Consequently, it is unlikely to be the ancestor of the modern horse; instead, it is a likely candidate for the ancestor of Astrohippus.
Dinohippus
Dinohippus was the most common species of Equidae in North America during the late Pliocene. It was originally thought to be monodactyl, but a 1981 fossil find in Nebraska shows some were tridactyl.
Plesippus
Plesippus is often considered an intermediate stage between Dinohippus and the extant genus, Equus.
The famous fossils found near Hagerman, Idaho, were originally thought to be a part of the genus Plesippus. Hagerman Fossil Beds (Idaho) is a Pliocene site, dating to about 3.5 mya. The fossilized remains were originally called Plesippus shoshonensis, but further study by paleontologists determined the fossils represented the oldest remains of the genus Equus. Their estimated average weight was 425 kg, roughly the size of an Arabian horse.
At the end of the Pliocene, the climate in North America began to cool significantly and most of the animals were forced to move south. One population of Plesippus moved across the Bering land bridge into Eurasia around 2.5 mya.
Modern horses
Equus
The genus Equus, which includes all extant equines, is believed to have evolved from Dinohippus, via the intermediate form Plesippus. One of the oldest species is Equus simplicidens, described as zebra-like with a donkey-shaped head. The oldest fossil to date is ~3.5 million years old, discovered in Idaho. The genus appears to have spread quickly into the Old World, with the similarly aged Equus livenzovensis documented from western Europe and Russia.
Molecular phylogenies indicate the most recent common ancestor of all modern equids (members of the genus Equus) lived ~5.6 (3.9–7.8) mya. Direct paleogenomic sequencing of a 700,000-year-old middle Pleistocene horse metapodial bone from Canada implies a more recent 4.07 Myr before present date for the most recent common ancestor (MRCA) within the range of 4.0 to 4.5 Myr BP. The oldest divergencies are the Asian hemiones (subgenus E. (Asinus), including the kulan, onager, and kiang), followed by the African zebras (subgenera E. (Dolichohippus), and E. (Hippotigris)). All other modern forms including the domesticated horse (and many fossil Pliocene and Pleistocene forms) belong to the subgenus E. (Equus) which diverged ~4.8 (3.2–6.5) million years ago.
Pleistocene horse fossils have been assigned to a multitude of species, with over 50 species of equines described from the Pleistocene of North America alone, although the taxonomic validity of most of these has been called into question. Recent genetic work on fossils has found evidence for only three genetically divergent equid lineages in Pleistocene North and South America. These results suggest all North American fossils of caballine-type horses (which also include the domesticated horse and Przewalski's horse of Europe and Asia), as well as South American fossils traditionally placed in the subgenus E. (Amerhippus) belong to the same species: E. ferus. Remains attributed to a variety of species and lumped as New World stilt-legged horses (including Haringtonhippus, E. tau, E. quinni and potentially North American Pleistocene fossils previously attributed to E. cf. hemiones, and E. (Asinus) cf. kiang) probably all belong to a second species endemic to North America, which despite a superficial resemblance to species in the subgenus E. (Asinus) (and hence occasionally referred to as North American ass) is closely related to E. ferus. Surprisingly, the third species, endemic to South America and traditionally referred to as Hippidion, originally believed to be descended from Pliohippus, was shown to be a third species in the genus Equus, closely related to the New World stilt-legged horse. The temporal and regional variation in body size and morphological features within each lineage indicates extraordinary intraspecific plasticity. Such environment-driven adaptative changes would explain why the taxonomic diversity of Pleistocene equids has been overestimated on morphoanatomical grounds.
According to these results, it appears the genus Equus evolved from a Dinohippus-like ancestor ~4–7 mya. It rapidly spread into the Old World and there diversified into the various species of asses and zebras. A North American lineage of the subgenus E. (Equus) evolved into the New World stilt-legged horse (NWSLH). Subsequently, populations of this species entered South America as part of the Great American Interchange shortly after the formation of the Isthmus of Panama, and evolved into the form currently referred to as Hippidion ~2.5 million years ago. Hippidion is thus only distantly related to the morphologically similar Pliohippus, which presumably became extinct during the Miocene. Both the NWSLH and Hippidium show adaptations to dry, barren ground, whereas the shortened legs of Hippidion may have been a response to sloped terrain. In contrast, the geographic origin of the closely related modern E. ferus is not resolved. However, genetic results on extant and fossil material of Pleistocene age indicate two clades, potentially subspecies, one of which had a holarctic distribution spanning from Europe through Asia and across North America and would become the founding stock of the modern domesticated horse. The other population appears to have been restricted to North America. However, one or more North American populations of E. ferus entered South America ~1.0–1.5 million years ago, leading to the forms currently known as E. (Amerhippus), which represent an extinct geographic variant or race of E. ferus.
Genome sequencing
Early sequencing studies of DNA revealed several genetic characteristics of Przewalski's horse that differ from what is seen in modern domestic horses, indicating neither is ancestor of the other, and supporting the status of Przewalski horses as a remnant wild population not derived from domestic horses. The evolutionary divergence of the two populations was estimated to have occurred about 45,000 YBP, while the archaeological record places the first horse domestication about 5,500 YBP by the ancient central-Asian Botai culture. The two lineages thus split well before domestication, probably due to climate, topography, or other environmental changes.
Several subsequent DNA studies produced partially contradictory results. A 2009 molecular analysis using ancient DNA recovered from archaeological sites placed Przewalski's horse in the middle of the domesticated horses, but a 2011 mitochondrial DNA analysis suggested that Przewalski's and modern domestic horses diverged some 160,000 years ago. An analysis based on whole genome sequencing and calibration with DNA from old horse bones gave a divergence date of 38–72 thousand years ago.
In June 2013, a group of researchers announced that they had sequenced the DNA of a 560–780 thousand year old horse, using material extracted from a leg bone found buried in permafrost in Canada's Yukon territory. Before this publication, the oldest nuclear genome that had been successfully sequenced was dated at 110–130 thousand years ago. For comparison, the researchers also sequenced the genomes of a 43,000-year-old Pleistocene horse, a Przewalski's horse, five modern horse breeds, and a donkey. Analysis of differences between these genomes indicated that the last common ancestor of modern horses, donkeys, and zebras existed 4 to 4.5 million years ago. The results also indicated that Przewalski's horse diverged from other modern types of horse about 43,000 years ago, and had never in its evolutionary history been domesticated.
A new analysis in 2018 involved genomic sequencing of ancient DNA from mid-fourth-millennium B.C.E. Botai domestic horses, as well as domestic horses from more recent archaeological sites, and comparison of these genomes with those of modern domestic and Przewalski's horses. The study revealed that Przewalski's horses not only belong to the same genetic lineage as those from the Botai culture, but were the feral descendants of these ancient domestic animals, rather than representing a surviving population of never-domesticated horses. The Botai horses were found to have made only negligible genetic contribution to any of the other ancient or modern domestic horses studied, which must then have arisen from an independent domestication involving a different wild horse population.
The karyotype of Przewalski's horse differs from that of the domestic horse by an extra chromosome pair because of the fission of domestic horse chromosome 5 to produce the Przewalski's horse chromosomes 23 and 24. In comparison, the chromosomal differences between domestic horses and zebras include numerous translocations, fusions, inversions and centromere repositioning. This gives Przewalski's horse the highest diploid chromosome number among all equine species. They can interbreed with the domestic horse and produce fertile offspring (65 chromosomes).
Pleistocene extinctions
Digs in western Canada have unearthed clear evidence horses existed in North America until about 12,000 years ago.
However, all Equidae in North America ultimately became extinct. The causes of this extinction (simultaneous with the extinctions of a variety of other American megafauna) have been a matter of debate. Given the suddenness of the event and because these mammals had been flourishing for millions of years previously, something quite unusual must have happened. The first main hypothesis attributes extinction to climate change. For example, in Alaska, beginning approximately 12,500 years ago, the grasses characteristic of a steppe ecosystem gave way to shrub tundra, which was covered with unpalatable plants. The other hypothesis suggests extinction was linked to overexploitation by newly arrived humans of naive prey that were not habituated to their hunting methods. The extinctions were roughly simultaneous with the end of the most recent glacial advance and the appearance of the big game-hunting Clovis culture. Several studies have indicated humans probably arrived in Alaska before or shortly before the local extinction of horses. However, it has been proposed that the steppe–tundra vegetation transition in Beringia may have been a consequence, rather than a cause, of the extinction of megafaunal grazers.
In Eurasia, horse fossils began occurring frequently again in archaeological sites in Kazakhstan and the southern Ukraine about 6,000 years ago. From then on, domesticated horses, as well as the knowledge of capturing, taming, and rearing horses, probably spread relatively quickly, with wild mares from several wild populations being incorporated en route.
Return to the Americas
Horses only returned to the Americas with Christopher Columbus in 1493. These were Iberian horses first brought to Hispaniola and later to Panama, Mexico, Brazil, Peru, Argentina, and, in 1538, Florida. The first horses to return to the main continent were 16 specifically identified horses brought by Hernán Cortés. Subsequent explorers, such as Coronado and De Soto, brought ever-larger numbers, some from Spain and others from breeding establishments set up by the Spanish in the Caribbean. Later, as Spanish missions were founded on the mainland, horses would eventually be lost or stolen, and proliferated into large herds of feral horses that became known as mustangs.
Details
Toes
The ancestors of the horse came to walk only on the end of the third toe and both side (second and fourth) "toes". Skeletal remnants show obvious wear on the back of both sides of metacarpal and metatarsal bones, commonly called the "splint bones". They are the remnants of the second and the fourth toes. Modern horses retain the splint bones; they are often believed to be useless attachments, but they in fact play an important role in supporting the carpal joints (front knees) and even the tarsal joints (hocks).
A 2018 study has found remnants of the remaining digits in the horse's hoof, suggesting a retention of all five digits (albeit in a "hourglass" arrangement where metacarpals/tarsals are present proximally and phalanges distally).
Teeth
Throughout the phylogenetic development, the teeth of the horse underwent significant changes. The type of the original omnivorous teeth with short, "bumpy" molars, with which the prime members of the evolutionary line distinguished themselves, gradually changed into the teeth common to herbivorous mammals. They became long (as much as 100 mm), roughly cubical molars equipped with flat grinding surfaces. In conjunction with the teeth, during the horse's evolution, the elongation of the facial part of the skull is apparent, and can also be observed in the backward-set eyeholes. In addition, the relatively short neck of the equine ancestors became longer, with equal elongation of the legs. Finally, the size of the body grew as well.
Coat color
The ancestral coat color of E. ferus was possibly a uniform dun, consistent with modern populations of Przewalski's horses. Pre-domestication variants including black and spotted have been inferred from cave wall paintings and confirmed by genomic analysis.
Domestication may have also led to more varieties of coat colors.
| Biology and health sciences | Equidae | Animals |
249930 | https://en.wikipedia.org/wiki/Heart%20failure | Heart failure | Heart failure (HF), also known as congestive heart failure (CHF), is a syndrome caused by an impairment in the heart's ability to fill with and pump blood.
Although symptoms vary based on which side of the heart is affected, HF typically presents with shortness of breath, excessive fatigue, and bilateral leg swelling. The severity of the heart failure is mainly decided based on ejection fraction and also measured by the severity of symptoms. Other conditions that have symptoms similar to heart failure include obesity, kidney failure, liver disease, anemia, and thyroid disease.
Common causes of heart failure include coronary artery disease, heart attack, high blood pressure, atrial fibrillation, valvular heart disease, excessive alcohol consumption, infection, and cardiomyopathy. These cause heart failure by altering the structure or the function of the heart or in some cases both. There are different types of heart failure: right-sided heart failure, which affects the right heart, left-sided heart failure, which affects the left heart, and biventricular heart failure, which affects both sides of the heart. Left-sided heart failure may be present with a reduced ejection fraction or with a preserved ejection fraction. Heart failure is not the same as cardiac arrest, in which blood flow stops completely due to the failure of the heart to pump.
Diagnosis is based on symptoms, physical findings, and echocardiography. Blood tests, and a chest x-ray may be useful to determine the underlying cause. Treatment depends on severity and case. For people with chronic, stable, or mild heart failure, treatment usually consists of lifestyle changes, such as not smoking, physical exercise, and dietary changes, as well as medications. In heart failure due to left ventricular dysfunction, angiotensin-converting-enzyme inhibitors, angiotensin II receptor blockers (ARBs), or angiotensin receptor-neprilysin inhibitors, along with beta blockers, mineralocorticoid receptor antagonists and SGLT2 inhibitors are recommended. Diuretics may also be prescribed to prevent fluid retention and the resulting shortness of breath. Depending on the case, an implanted device such as a pacemaker or implantable cardiac defibrillator may sometimes be recommended. In some moderate or more severe cases, cardiac resynchronization therapy (CRT) or cardiac contractility modulation may be beneficial. In severe disease that persists despite all other measures, a cardiac assist device ventricular assist device, or, occasionally, heart transplantation may be recommended.
Heart failure is a common, costly, and potentially fatal condition, and is the leading cause of hospitalization and readmission in older adults. Heart failure often leads to more drastic health impairments than the failure of other, similarly complex organs such as the kidneys or liver. In 2015, it affected about 40 million people worldwide. Overall, heart failure affects about 2% of adults, and more than 10% of those over the age of 70. Rates are predicted to increase.
The risk of death in the first year after diagnosis is about 35%, while the risk of death in the second year is less than 10% in those still alive. The risk of death is comparable to that of some cancers. In the United Kingdom, the disease is the reason for 5% of emergency hospital admissions. Heart failure has been known since ancient times in Egypt; it is mentioned in the Ebers Papyrus around 1550 BCE.
Definition
When the heart functions poorly as a pump and does not circulate blood adequately via the circulatory system to meet the demands of the body the term cardiovascular insufficiency is sometimes used. This generally leads to the syndrome of heart failure, a combination of signs and symptoms when the heart functions poorly as a pump. This leads to high filling pressure in the left atrium with fluid accumulation and water retention. Most of the visible signs of heart failure are the result of this fluid accumulation (edema) and the adjective congestion is added to the definition of heart failure. Impaired ejection can lead to inadequate blood flow to the body tissues, resulting in ischemia.
Signs and symptoms
Congestive heart failure is a pathophysiological condition in which the heart's output is insufficient to meet the needs of the body and lungs. The term "congestive heart failure" is often used because one of the most common symptoms is congestion or fluid accumulation in the tissues and veins of the lungs or other parts of a person's body.
Congestion manifests itself particularly in the form of fluid accumulation and swelling (edema), in the form of peripheral edema (causing swollen limbs and feet) and pulmonary edema (causing difficulty breathing) and ascites (swollen abdomen). Pulse pressure, which is the difference between the systolic ("top number") and diastolic ("bottom number") blood pressures, is often low/narrow (i.e. 25% or less of the level of the systolic) in people with heart failure, and this can be an early warning sign.
Symptoms of heart failure are traditionally divided into left-sided and right-sided because the left and right ventricles supply different parts of the circulation. In biventricular heart failure, both sides of the heart are affected. Left-sided heart failure is the more common.
Left-sided failure
The left side of the heart takes oxygen-rich blood from the lungs and pumps it to the rest of the circulatory system in the body (except for the pulmonary circulation). Failure of the left side of the heart causes blood to back up into the lungs, causing breathing difficulties and fatigue due to an insufficient supply of oxygenated blood. Common respiratory signs include increased respiratory rate and labored breathing (nonspecific signs of shortness of breath). Rales or crackles are heard initially in the lung bases and when severe in all lung fields indicate the development of pulmonary edema (fluid in the alveoli). Cyanosis, indicates deficiency of oxygen in the blood, is a late sign of extremely severe pulmonary edema.
Other signs of left ventricular failure include a laterally displaced apex beat (which occurs when the heart is enlarged) and a gallop rhythm (additional heart sounds), which may be heard as a sign of increased blood flow or increased intracardiac pressure. Heart murmurs may indicate the presence of valvular heart disease, either as a cause (e.g., aortic stenosis) or as a consequence (e.g., mitral regurgitation) of heart failure.
Reverse insufficiency of the left ventricle causes congestion in the blood vessels of the lungs so that symptoms are predominantly respiratory. Reverse insufficiency can be divided into the failure of the left atrium, the left ventricle, or both within the left circuit. Patients will experience shortness of breath (dyspnea) on exertion and, in severe cases, dyspnea at rest. Increasing breathlessness while lying down, called orthopnea, also occurs. It can be measured by the number of pillows required to lie comfortably, with extreme cases of orthopnea forcing the patient to sleep sitting up. Another symptom of heart failure is paroxysmal nocturnal dyspnea: a sudden nocturnal attack of severe shortness of breath, usually occurring several hours after falling asleep. There may be "cardiac asthma" or wheezing. Impaired left ventricular forward function can lead to symptoms of poor systemic perfusion such as dizziness, confusion, and cool extremities at rest. Loss of consciousness may also occur due to loss of blood supply to the brain.
Right-sided failure
Right-sided heart failure is often caused by pulmonary heart disease (cor pulmonale), which is typically caused by issues with pulmonary circulation such as pulmonary hypertension or pulmonic stenosis. Physical examination may reveal pitting peripheral edema, ascites, liver enlargement, and spleen enlargement. Jugular venous pressure is frequently assessed as a marker of fluid status, which can be accentuated by testing hepatojugular reflux. If the right ventricular pressure is increased, a parasternal heave which causes the compensatory increase in contraction strength may be present.
Backward failure of the right ventricle leads to congestion of systemic capillaries. This generates excess fluid accumulation in the body. This causes swelling under the skin (peripheral edema or anasarca) and usually affects the dependent parts of the body first, causing foot and ankle swelling in people who are standing up and sacral edema in people who are predominantly lying down. Nocturia (frequent night-time urination) may occur when fluid from the legs is returned to the bloodstream while lying down at night. In progressively severe cases, ascites (fluid accumulation in the abdominal cavity causing swelling) and liver enlargement may develop. Significant liver congestion may result in impaired liver function (congestive hepatopathy), jaundice, and coagulopathy (problems of decreased or increased blood clotting).
Biventricular failure
Dullness of the lung fields when percussed and reduced breath sounds at the base of the lungs may suggest the development of a pleural effusion (fluid collection between the lung and the chest wall). Though it can occur in isolated left- or right-sided heart failure, it is more common in biventricular failure because pleural veins drain into both the systemic and pulmonary venous systems. When unilateral, effusions are often right-sided.
If a person with a failure of one ventricle lives long enough, it will tend to progress to failure of both ventricles. For example, left ventricular failure allows pulmonary edema and pulmonary hypertension to occur, which increases stress on the right ventricle. Though still harmful, right ventricular failure is not as deleterious to the left side.
Causes
Since heart failure is a syndrome and not a disease, establishing the underlying cause is vital to diagnosis and treatment. In heart failure, the structure or the function of the heart or in some cases both are altered. Heart failure is the potential end stage of all heart diseases.
Common causes of heart failure include coronary artery disease, including a previous myocardial infarction (heart attack), high blood pressure, atrial fibrillation, valvular heart disease, excess alcohol use, infection, and cardiomyopathy of an unknown cause. In addition, viral infection and subsequent inflammation of the heart's myocardial tissue (termed myocarditis) can similarly contribute to the development of heart failure. Genetic predisposition plays an important role. If more than one cause is present, progression is more likely and prognosis is worse.
Heart damage can predispose a person to develop heart failure later in life and has many causes including systemic viral infections (e.g., HIV), chemotherapeutic agents such as daunorubicin, cyclophosphamide, trastuzumab and substance use disorders of substances such as alcohol, cocaine, and methamphetamine. An uncommon cause is exposure to certain toxins such as lead and cobalt. Additionally, infiltrative disorders such as amyloidosis and connective tissue diseases such as systemic lupus erythematosus have similar consequences. Obstructive sleep apnea (a condition of sleep wherein disordered breathing overlaps with obesity, hypertension, and/or diabetes) is regarded as an independent cause of heart failure. Recent reports from clinical trials have also linked variation in blood pressure to heart failure and cardiac changes that may give rise to heart failure.
High-output heart failure
High-output heart failure happens when the amount of blood pumped out is more than typical and the heart cannot keep up. This can occur in overload situations such as blood or serum infusions, kidney diseases, chronic severe anemia, beriberi (vitamin B1/thiamine deficiency), hyperthyroidism, cirrhosis, Paget's disease, multiple myeloma, arteriovenous fistulae, or arteriovenous malformations.
Acute decompensation
Chronic stable heart failure may easily decompensate (fail to meet the body's metabolic needs). This most commonly results from a concurrent illness (such as myocardial infarction (a heart attack) or pneumonia), abnormal heart rhythms, uncontrolled hypertension, or a person's failure to maintain a fluid restriction, diet, or medication.
Other factors that may worsen CHF include: anemia, hyperthyroidism, excessive fluid or salt intake, and medication such as NSAIDs and thiazolidinediones. NSAIDs increase the risk twofold.
Medications
A number of medications may cause or worsen the disease. This includes NSAIDs, COX-2 inhibitors, a number of anesthetic agents such as ketamine, thiazolidinediones, some cancer medications, several antiarrhythmic medications, pregabalin, alpha-2 adrenergic receptor agonists, minoxidil, itraconazole, cilostazol, anagrelide, stimulants (e.g., methylphenidate), tricyclic antidepressants, lithium, antipsychotics, dopamine agonists, TNF inhibitors, calcium channel blockers (especially verapamil and diltiazem), salbutamol, and tamsulosin.
By inhibiting the formation of prostaglandins, NSAIDs may exacerbate heart failure through several mechanisms, including promotion of fluid retention, increasing blood pressure, and decreasing a person's response to diuretic medications. Similarly, the ACC/AHA recommends against using COX-2 inhibitor medications in people with heart failure. Thiazolidinediones have been strongly linked to new cases of heart failure and worsening of pre-existing congestive heart failure due to their association with weight gain and fluid retention. Certain calcium channel blockers, such as diltiazem and verapamil, are known to decrease the force with which the heart ejects blood, thus are not recommended in people with heart failure with a reduced ejection fraction.
Breast cancer patients are at high risk of heart failure due to several factors. After analyzing data from 26 studies (836,301 patients), the recent meta-analysis found that breast cancer survivors demonstrated a higher risk heart failure within first ten years after diagnosis (hazard ratio = 1.21; 95% CI: 1.1, 1.33). The pooled incidence of heart failure in breast cancer survivors was 4.44 (95% CI 3.33-5.92) per 1000 person-years of follow-up.
Supplements
Certain alternative medicines carry a risk of exacerbating existing heart failure, and are not recommended. This includes aconite, ginseng, gossypol, gynura, licorice, lily of the valley, tetrandrine, and yohimbine. Aconite can cause abnormally slow heart rates and abnormal heart rhythms such as ventricular tachycardia. Ginseng can cause abnormally low or high blood pressure and may interfere with the effects of diuretic medications. Gossypol can increase the effects of diuretics, leading to toxicity.
Gynura can cause low blood pressure. Licorice can worsen heart failure by increasing blood pressure and promoting fluid retention. Lily of the Valley can cause abnormally slow heart rates with mechanisms similar to those of digoxin. Tetrandrine can lower blood pressure by inhibiting L-type calcium channels. Yohimbine can exacerbate heart failure by increasing blood pressure through alpha-2 adrenergic receptor antagonism.
Pathophysiology
Heart failure is caused by any condition that reduces the efficiency of the heart muscle, through damage or overloading. Over time, these increases in workload, which are mediated by long-term activation of neurohormonal systems such as the renin–angiotensin system and the sympathoadrenal system, lead to fibrosis, dilation, and structural changes in the shape of the left ventricle from elliptical to spherical.
The heart of a person with heart failure may have a reduced force of contraction due to overloading of the ventricle. In a normal heart, increased filling of the ventricle results in increased contraction force by the Frank–Starling law of the heart, and thus a rise in cardiac output. In heart failure, this mechanism fails, as the ventricle is loaded with blood to the point where heart muscle contraction becomes less efficient. This is due to the reduced ability to cross-link actin and myosin myofilaments in over-stretched heart muscle.
Diagnosis
No diagnostic criteria have been agreed on as the gold standard for heart failure, especially heart failure with preserved ejection fraction (HFpEF).
In the UK, the National Institute for Health and Care Excellence recommends measuring N-terminal pro-BNP (NT-proBNP) followed by an ultrasound of the heart if positive. In Europe, the European Society of Cardiology, and in the United States, the AHA/ACC/HFSA, recommend measuring NT-proBNP or BNP followed by an ultrasound of the heart if positive. This is recommended in those with symptoms consistent with heart failure such as shortness of breath.
The European Society of Cardiology defines the diagnosis of heart failure as symptoms and signs consistent with heart failure in combination with "objective evidence of cardiac structural or functional abnormalities". This definition is consistent with an international 2021 report termed "Universal Definition of Heart Failure". Score-based algorithms have been developed to help in the diagnosis of HFpEF, which can be challenging for physicians to diagnose. The AHA/ACC/HFSA defines heart failure as symptoms and signs consistent with heart failure in combination with shown "structural and functional alterations of the heart as the underlying cause for the clinical presentation", for HFmrEF and HFpEF specifically requiring "evidence of spontaneous or provokable increased left ventricle filling pressures".
Algorithms
The European Society of Cardiology has developed a diagnostic algorithm for HFpEF, named HFA-PEFF. HFA-PEFF considers symptoms and signs, typical clinical demographics (obesity, hypertension, diabetes, elderly, atrial fibrillation), and diagnostic laboratory tests, ECG, and echocardiography.
Classification
"Left", "right" and mixed heart failure
One historical method of categorizing heart failure is by the side of the heart involved (left heart failure versus right heart failure). Right heart failure was thought to compromise blood flow to the lungs compared to left heart failure compromising blood flow to the aorta and consequently to the brain and the remainder of the body's systemic circulation. However, mixed presentations are common, and left heart failure is a common cause of right heart failure.
By ejection fraction
A more accurate classification of heart failure type is made by measuring ejection fraction, or the proportion of blood pumped out of the heart during a single contraction. Ejection fraction is given as a percentage with the normal range being between 50 and 75%. The types are:
Heart failure with reduced ejection fraction (HFrEF): Synonyms no longer recommended are "heart failure due to left ventricular systolic dysfunction" and "systolic heart failure". HFrEF is associated with an ejection fraction less than 40%.
Heart failure with mildly reduced ejection fraction (HFmrEF), previously called "heart failure with mid-range ejection fraction", is defined by an ejection fraction of 41–49%.
Heart failure with preserved ejection fraction (HFpEF): Synonyms no longer recommended include "diastolic heart failure" and "heart failure with normal ejection fraction." HFpEF occurs when the left ventricle contracts normally during systole, but the ventricle is stiff and does not relax normally during diastole, which impairs filling.
Heart failure with recovered ejection fraction (HFrecovEF or HFrecEF): patients previously with HFrEF with complete normalization of left ventricular ejection (≥50%).
Heart failure may also be classified as acute or chronic. Chronic heart failure is a long-term condition, usually kept stable by the treatment of symptoms. Acute decompensated heart failure is a worsening of chronic heart failure symptoms, which can result in acute respiratory distress. High-output heart failure can occur when there is increased cardiac demand that results in increased left ventricular diastolic pressure which can develop into pulmonary congestion (pulmonary edema).
Several terms are closely related to heart failure and may be the cause of heart failure, but should not be confused with it. Cardiac arrest and asystole refer to situations in which no cardiac output occurs at all. Without urgent treatment, these events result in sudden death. Myocardial infarction ("Heart attack") refers to heart muscle damage due to insufficient blood supply, usually as a result of a blocked coronary artery. Cardiomyopathy refers specifically to problems within the heart muscle, and these problems can result in heart failure. Ischemic cardiomyopathy implies that the cause of muscle damage is coronary artery disease. Dilated cardiomyopathy implies that the muscle damage has resulted in enlargement of the heart. Hypertrophic cardiomyopathy involves enlargement and thickening of the heart muscle.
Ultrasound
An echocardiogram (ultrasound of the heart) is commonly used to support a clinical diagnosis of heart failure. This can determine the stroke volume (SV, the amount of blood in the heart that exits the ventricles with each beat), the end-diastolic volume (EDV, the total amount of blood at the end of diastole), and the SV in proportion to the EDV, a value known as the ejection fraction (EF). In pediatrics, the shortening fraction is the preferred measure of systolic function. Normally, the EF should be between 50 and 70%; in systolic heart failure, it drops below 40%. Echocardiography can also identify valvular heart disease and assess the state of the pericardium (the connective tissue sac surrounding the heart). Echocardiography may also aid in deciding specific treatments, such as medication, insertion of an implantable cardioverter-defibrillator, or cardiac resynchronization therapy. Echocardiography can also help determine if acute myocardial ischemia is the precipitating cause, and may manifest as regional wall motion abnormalities on echo.
Chest X-ray
Chest X-rays are frequently used to aid in the diagnosis of CHF. In a person who is compensated, this may show cardiomegaly (visible enlargement of the heart), quantified as the cardiothoracic ratio (proportion of the heart size to the chest). In left ventricular failure, evidence may exist of vascular redistribution (upper lobe blood diversion or cephalization), Kerley lines, cuffing of the areas around the bronchi, and interstitial edema. Ultrasound of the lung may also detect Kerley lines.
Electrophysiology
An electrocardiogram (ECG or EKG) may be used to identify arrhythmias, ischemic heart disease, right and left ventricular hypertrophy, and presence of conduction delay or abnormalities (e.g. left bundle branch block). Although these findings are not specific to the diagnosis of heart failure, a normal ECG virtually excludes left ventricular systolic dysfunction.
Blood tests
N-terminal pro-BNP (NT-proBNP) is the favored biomarker for the diagnosis of heart failure, according to guidelines published 2018 by NICE in the UK. Brain natriuretic peptide 32 (BNP) is another biomarker commonly tested for heart failure. An elevated NT-proBNP or BNP is a specific test indicative of heart failure. Additionally, NT-proBNP or BNP can be used to differentiate between causes of dyspnea due to heart failure from other causes of dyspnea. If a myocardial infarction is suspected, various cardiac markers may be used.
Blood tests routinely performed include electrolytes (sodium, potassium), measures of kidney function, liver function tests, thyroid function tests, a complete blood count, and often C-reactive protein if infection is suspected.
Hyponatremia (low serum sodium concentration) is common in heart failure. Vasopressin levels are usually increased, along with renin, angiotensin II, and catecholamines to compensate for reduced circulating volume due to inadequate cardiac output. This leads to increased fluid and sodium retention in the body; the rate of fluid retention is higher than the rate of sodium retention in the body, this phenomenon causes hypervolemic hyponatremia (low sodium concentration due to high body fluid retention). This phenomenon is more common in older women with low body mass. Severe hyponatremia can result in accumulation of fluid in the brain, causing cerebral edema and intracranial hemorrhage.
Angiography
Angiography is the X-ray imaging of blood vessels, which is done by injecting contrast agents into the bloodstream through a thin plastic tube (catheter), which is placed directly in the blood vessel. X-ray images are called angiograms. Heart failure may be the result of coronary artery disease, and its prognosis depends in part on the ability of the coronary arteries to supply blood to the myocardium (heart muscle). As a result, coronary catheterization may be used to identify possibilities for revascularisation through percutaneous coronary intervention or bypass surgery.
Staging
Heart failure is commonly stratified by the degree of functional impairment conferred by the severity of the heart failure, as reflected in the New York Heart Association (NYHA) functional classification. The NYHA functional classes (I–IV) begin with class I, which is defined as a person who experiences no limitation in any activities and has no symptoms from ordinary activities. People with NYHA class II heart failure have slight, mild limitations with everyday activities; the person is comfortable at rest or with mild exertion. With NYHA class III heart failure, a marked limitation occurs with any activity; the person is comfortable only at rest. A person with NYHA class IV heart failure is symptomatic at rest and becomes quite uncomfortable with any physical activity. This score documents the severity of symptoms and can be used to assess response to treatment. While its use is widespread, the NYHA score is not very reproducible and does not reliably predict walking distance or exercise tolerance on formal testing.
In its 2001 guidelines, the American College of Cardiology/American Heart Association working group introduced four stages of heart failure:
Stage A: People at high risk for developing HF in the future, but no functional or structural heart disorder
Stage B: A structural heart disorder, but no symptoms at any stage
Stage C: Previous or current symptoms of heart failure in the context of an underlying structural heart problem, but managed with medical treatment
Stage D: Advanced disease requiring hospital-based support, a heart transplant, or palliative care
The ACC staging system is useful since stage A encompasses "pre-heart failure" – a stage where intervention with treatment can presumably prevent progression to overt symptoms. ACC stage A does not have a corresponding NYHA class. ACC stage B would correspond to NYHA class I. ACC stage C corresponds to NYHA class II and III, while ACC stage D overlaps with NYHA class IV.
The degree of coexisting illness: i.e. heart failure/systemic hypertension, heart failure/pulmonary hypertension, heart failure/diabetes, heart failure/kidney failure, etc.
Whether the problem is primarily increased venous back pressure (preload), or failure to supply adequate arterial perfusion (afterload)
Whether the abnormality is due to low cardiac output with high systemic vascular resistance or high cardiac output with low vascular resistance (low-output heart failure vs. high-output heart failure)
Histopathology
Histopathology can diagnose heart failure in autopsies. The presence of siderophages indicates chronic left-sided heart failure, but is not specific for it. It is also indicated by congestion of the pulmonary circulation.
Prevention
A person's risk of developing heart failure is inversely related to the level of physical activity. Those who achieved at least 500 MET-minutes/week (the recommended minimum by U.S. guidelines) had lower heart failure risk than individuals who did not report exercising during their free time; the reduction in heart failure risk was even greater in those who engaged in higher levels of physical activity than the recommended minimum.
Heart failure can also be prevented by lowering high blood pressure and high blood cholesterol, and by controlling diabetes. Maintaining a healthy weight, and decreasing sodium, alcohol, and sugar intake, may help. Additionally, avoiding tobacco use has been shown to lower the risk of heart failure.
According to Johns Hopkins and the American Heart Association there are a few ways to help prevent a cardiac event. Johns Hopkins states that stopping tobacco use, reducing high blood pressure, physical activity, and nutrition can drastically affect the chances of developing heart disease. High blood pressure accounts for most cardiovascular deaths. High blood pressure can be lowered into the normal range by making dietary decisions such as consuming less salt. Exercise also helps to bring blood pressure back down. One of the best ways to help avoid heart failure is to promote healthier eating habits like eating more vegetables, fruits, grains, and lean protein.
Diabetes is a major risk factor for heart failure. For women with Coronary Heart disease (CHD), diabetes was the strongest risk factor for heart failure. Diabetic women with depressed creatinine clearance or elevated BMI were at the highest risk of heart failure. While the annual incidence rate of heart failure for non-diabetic women with no risk factors is 0.4%, the annual incidence rate for diabetic women with elevated body mass index (BMI) and depressed creatinine clearance was 7% and 13%, respectively.
Management
Treatment focuses on improving the symptoms and preventing the progression of the disease. Reversible causes of heart failure also need to be addressed (e.g. infection, alcohol ingestion, anemia, thyrotoxicosis, arrhythmia, and hypertension). Treatments include lifestyle and pharmacological modalities, and occasionally various forms of device therapy. Rarely, cardiac transplantation is used as an effective treatment when heart failure has reached the end stage.
Acute decompensation
In acute decompensated heart failure, the immediate goal is to re-establish adequate perfusion and oxygen delivery to end organs. This entails ensuring that airway, breathing, and circulation are adequate. Immediate treatments usually involve some combination of vasodilators such as nitroglycerin, diuretics such as furosemide, and possibly noninvasive positive pressure ventilation. Supplemental oxygen is indicated in those with oxygen saturation levels below 90%, but is not recommended in those with normal oxygen levels in the normal atmosphere.
Chronic management
The goals of treatment for people with chronic heart failure are prolonging life, preventing acute decompensation, and reducing symptoms, allowing for greater activity.
Heart failure can result from a variety of conditions. In considering therapeutic options, excluding reversible causes is of primary importance, including thyroid disease, anemia, chronic tachycardia, alcohol use disorder, hypertension, and dysfunction of one or more heart valves. Treatment of the underlying cause is usually the first approach to treating heart failure. In most cases, though, either no primary cause is found or treatment of the primary cause does not restore normal heart function. In these cases, behavioral, medical and device treatment strategies exist that can provide a significant improvement in outcomes, including the relief of symptoms, exercise tolerance, and a decrease in the likelihood of hospitalization or death. Breathlessness rehabilitation for chronic obstructive pulmonary disease and heart failure has been proposed with exercise training as a core component. Rehabilitation should also include other interventions to address shortness of breath including the psychological and educational needs of people and the needs of caregivers. Iron supplementation appears to reduce hospitalization but not all-cause mortality in patients with iron deficiency and heart
failure.
Advance care planning
The latest evidence indicates that advance care planning (ACP) may help to increase documentation by medical staff regarding discussions with participants and improve an individual's depression. This involves discussing an individual's future care plan, preferences, and values. The findings are, however, based on low-quality evidence.
Monitoring
The various measures often used to assess the progress of people being treated for heart failure include fluid balance (calculation of fluid intake and excretion) and monitoring body weight (which in the shorter term reflects fluid shifts). Remote monitoring can be effective to reduce complications for people with heart failure.
Lifestyle
Behavior modification is a primary consideration in chronic heart failure management programs, with dietary guidelines regarding fluid and salt intake. Fluid restriction is important to reduce fluid retention in the body and to correct the hyponatremic status of the body. The evidence of the benefit of reducing salt, however, is poor as of 2018. Thirst is a common and burdensome symptom for patients to cope with. Chewing gum is an effective intervention to relieve thirst in patients experiencing heart failure, although patient acceptability remains an issue.
Exercise and physical activity
Exercise should be encouraged and tailored to suit an individual's capabilities. A meta-analysis found that center-based group interventions delivered by a physiotherapist help promote physical activity in HF. There is a need for additional training for physiotherapists in delivering behavior change intervention alongside an exercise program. An intervention is expected to be more efficacious in encouraging physical activity than the usual care if it includes Prompts and cues to walk or exercise, like a phone call or a text message. It is helpful if a trusted clinician provides explicit advice to engage in physical activity (Credible source). Another highly effective strategy is to place objects that will serve as a cue to engage in physical activity in the person's everyday environment (Adding object to the environment; e.g., exercise step or treadmill). Encouragement to walk or exercise in various settings beyond CR (e.g., home, neighborhood, parks) is also promising (Generalisation of target behavior). Additional promising strategies are Graded tasks (e.g., gradual increase in intensity and duration of exercise training), Self-monitoring, Monitoring of physical activity by others without feedback, Action planning, and Goal-setting. The inclusion of regular physical conditioning as part of a cardiac rehabilitation program can significantly improve quality of life and reduce the risk of hospital admission for worsening symptoms, but no evidence shows a reduction in mortality rates as a result of exercise.
Home visits and regular monitoring at heart-failure clinics reduce the need for hospitalization and improve life expectancy.
Medication
Quadruple medical therapy using a combination of angiotensin receptor-neprilysin inhibitors (ARNI), beta blockers, mineralocorticoid receptor antagonists (MRA), and sodium/glucose cotransporter 2 inhibitors (SGLT2 inhibitors) is the standard of care as of 2021 for heart failure with reduced ejection fraction (HFrEF).
There is no convincing evidence for pharmacological treatment of heart failure with preserved ejection fraction (HFpEF). Medication for HFpEF is symptomatic treatment with diuretics to treat congestion. Managing risk factors and comorbidities such as hypertension is recommended in HFpEF.
Inhibitors of the renin–angiotensin system (RAS) are recommended for heart failure. The angiotensin receptor-neprilysin inhibitors (ARNI) sacubitril/valsartan is recommended as the first choice of RAS inhibitors in American guidelines published by AHA/ACC in 2022. Use of ACE inhibitor, or angiotensin receptor blockers (ARBs) if the person develops a long-term cough as a side effect of the ACE-I, is associated with improved survival, fewer hospitalizations for heart failure exacerbations, and improved quality of life in people with heart failure. European guidelines published by ESC in 2021 recommends that ARNI should be used in those who still have symptoms while on an ACE-I or ARB, beta blocker, and a mineralocorticoid receptor antagonist. Use of the combination agent ARNI requires the cessation of ACE-I or ARB therapy at least 36 hours before its initiation.
Beta-adrenergic blocking agents (beta blockers) add to the improvement in symptoms and mortality provided by ACE-I/ARB. The mortality benefits of beta blockers in people with systolic dysfunction who also have atrial fibrillation is more limited than in those who do not have it. If the ejection fraction is not diminished (HFpEF), the benefits of beta blockers are more modest; a decrease in mortality has been observed, but reduction in hospital admission for uncontrolled symptoms has not been observed.
In people who are intolerant of ACE-I and ARB or who have significant kidney dysfunction, the use of combined hydralazine and a long-acting nitrate, such as isosorbide dinitrate, is an effective alternate strategy. This regimen has been shown to reduce mortality in people with moderate heart failure. It is especially beneficial in the black population.
Use of a mineralocorticoid antagonist, such as spironolactone or eplerenone, in addition to beta blockers and ACE-I, can improve symptoms and reduce mortality in people with symptomatic heart failure with reduced ejection fraction (HFrEF).
SGLT2 inhibitors are used for heart failure with reduced ejection fraction as they have demonstrated benefits in reducing hospitalizations and mortality, regardless of whether an individual has comorbid Type 2 Diabetes or not.
Other medications
Second-line medications for CHF do not confer a mortality benefit. Digoxin is one such medication. Its narrow therapeutic window, a high degree of toxicity, and the failure of multiple trials to show a mortality benefit have reduced its role in clinical practice. It is now used in only a small number of people with refractory symptoms, who are in atrial fibrillation, and/or who have chronic hypotension.
Diuretics have been a mainstay of treatment against symptoms of fluid accumulation, and include diuretics classes such as loop diuretics (such as furosemide), thiazide-like diuretics, and potassium-sparing diuretics. Although widely used, evidence on their efficacy and safety is limited, except for mineralocorticoid antagonists such as spironolactone.
Anemia is an independent factor in mortality in people with chronic heart failure. Treatment of anemia significantly improves the quality of life for those with heart failure, often with a reduction in severity of the NYHA classification, and also improves mortality rates. The European Society of Cardiology recommends screening for iron deficiency and treating with intravenous iron if deficiency is found.
The decision to anticoagulate people with HF, typically with left ventricular ejection fractions <35% is debated, but generally, people with coexisting atrial fibrillation, a prior embolic event, or conditions that increase the risk of an embolic event such as amyloidosis, left ventricular noncompaction, familial dilated cardiomyopathy, or a thromboembolic event in a first-degree relative.
Vasopressin receptor antagonists can also treat heart failure. Conivaptan is the first medication approved by the US Food and Drug Administration for the treatment of euvolemic hyponatremia in those with heart failure. In rare cases hypertonic 3% saline together with diuretics may be used to correct hyponatremia.
Ivabradine is recommended for people with symptomatic heart failure with reduced left ventricular ejection fraction who are receiving optimized guideline-directed therapy (as above) including the maximum tolerated dose of beta-blocker, have a normal heart rhythm and continue to have a resting heart rate above 70 beats per minute. Ivabradine has been found to reduce the risk of hospitalization for heart failure exacerbations in this subgroup of people with heart failure.
Implanted devices
In people with severe cardiomyopathy (left ventricular ejection fraction below 35%), or in those with recurrent VT or malignant arrhythmias, treatment with an automatic implantable cardioverter-defibrillator (AICD) is indicated to reduce the risk of severe life-threatening arrhythmias. The AICD does not improve symptoms or reduce the incidence of malignant arrhythmias but does reduce mortality from those arrhythmias, often in conjunction with antiarrhythmic medications. In people with left ventricular ejection (LVEF) below 35%, the incidence of ventricular tachycardia or sudden cardiac death is high enough to warrant AICD placement. Its use is therefore recommended in AHA/ACC guidelines.
Cardiac contractility modulation (CCM) is a treatment for people with moderate to severe left ventricular systolic heart failure (NYHA classes II–IV), which enhances both the strength of ventricular contraction and the heart's pumping capacity. The CCM mechanism is based on stimulation of the cardiac muscle by nonexcitatory electrical signals, which are delivered by a pacemaker-like device. CCM is particularly suitable for the treatment of heart failure with normal QRS complex duration (120 ms or less) and has been demonstrated to improve the symptoms, quality of life, and exercise tolerance. CCM is approved for use in Europe, and was approved by the Food and Drug Administration for use in the United States in 2019.
About one-third of people with an LVEF below 35% have markedly altered conduction to the ventricles, resulting in dyssynchronous depolarization of the right and left ventricles. This is especially problematic in people with left bundle branch block (blockage of one of the two primary conducting fiber bundles that originate at the base of the heart and carry depolarizing impulses to the left ventricle). Using a special pacing algorithm, biventricular cardiac resynchronization therapy (CRT) can initiate a normal sequence of ventricular depolarization. In people with LVEF below 35% and prolonged QRS duration on ECG (LBBB or QRS of 150 ms or more), an improvement in symptoms and mortality occurs when CRT is added to standard medical therapy. However, in the two-thirds of people without prolonged QRS duration, CRT may be harmful.
Surgical therapies
People with the most severe heart failure may be candidates for ventricular assist devices, which have commonly been used as a bridge to heart transplantation but have been used more recently as a destination treatment for advanced heart failure.
In select cases, heart transplantation can be considered. While this may resolve the problems associated with heart failure, the person must generally remain on an immunosuppressive regimen to prevent rejection, which has its own significant downsides. A major limitation of this treatment option is the scarcity of hearts available for transplantation.
Palliative care
People with heart failure often have significant symptoms, such as shortness of breath and chest pain. Palliative care should be initiated early in the HF trajectory, and should not be an option of last resort. Palliative care can not only provide symptom management, but also assist with advanced care planning, goals of care in the case of a significant decline, and making sure the person has a medical power of attorney and discussed his or her wishes with this individual. A 2016 and 2017 review found that palliative care is associated with improved outcomes, such as quality of life, symptom burden, and satisfaction with care.
Without transplantation, heart failure may not be reversible and heart function typically deteriorates with time. The growing number of people with stage IV heart failure (intractable symptoms of fatigue, shortness of breath, or chest pain at rest despite optimal medical therapy) should be considered for palliative care or hospice, according to American College of Cardiology/American Heart Association guidelines.
Prognosis
Prognosis in heart failure can be assessed in multiple ways, including clinical prediction rules and cardiopulmonary exercise testing. Clinical prediction rules use a composite of clinical factors such as laboratory tests and blood pressure to estimate prognosis. Among several clinical prediction rules for prognosticating acute heart failure, the 'EFFECT rule' slightly outperformed other rules in stratifying people and identifying those at low risk of death during hospitalization or within 30 days. Easy methods for identifying people that are low-risk are:
ADHERE Tree rule indicates that people with blood urea nitrogen < 43 mg/dL and systolic blood pressure at least 115 mm Hg have less than 10% chance of inpatient death or complications.
BWH rule indicates that people with systolic blood pressure over 90 mm Hg, respiratory rate of 30 or fewer breaths per minute, serum sodium over 135 mmol/L, and no new ST–T wave changes have less than 10% chance of inpatient death or complications.
A crucial method for assessing prognosis in people with advanced heart failure is cardiopulmonary exercise testing (CPX testing). CPX testing is usually required before heart transplantation as an indicator of prognosis. CPX testing involves the measurement of exhaled oxygen and carbon dioxide during exercise. The peak oxygen consumption (VO2 max) is used as an indicator of prognosis. As a general rule, a VO2 max less than 12–14 cc/kg/min indicates poor survival and suggests that the person may be a candidate for a heart transplant. People with a VO2 max <10 cc/kg/min have a poorer prognosis. The most recent International Society for Heart and Lung Transplantation guidelines also suggest two other parameters that can be used for evaluation of prognosis in advanced heart failure, the heart failure survival score and the use of a criterion of VE/VCO2 slope > 35 from the CPX test. The heart failure survival score is calculated using a combination of clinical predictors and the VO2 max from the CPX test.
Heart failure is associated with significantly reduced physical and mental health, resulting in a markedly decreased quality of life. With the exception of heart failure caused by reversible conditions, the condition usually worsens with time. Although some people survive many years, progressive disease is associated with an overall annual mortality rate of 10%.
Around 18 of every 1000 persons will experience an ischemic stroke during the first year after diagnosis of HF. As the duration of follow-up increases, the stroke rate rises to nearly 50 strokes per 1000 cases of HF by 5 years.
Epidemiology
In 2022, heart failure affected about 64 million people globally. Overall, around 2% of adults have heart failure. In those over the age of 75, rates are greater than 10%.
Rates are predicted to increase. Increasing rates are mostly because of increasing lifespan, but also because of increased risk factors (hypertension, diabetes, dyslipidemia, and obesity) and improved survival rates from other types of cardiovascular disease (myocardial infarction, valvular disease, and arrhythmias). Heart failure is the leading cause of hospitalization in people older than 65.
United States
In the United States, heart failure affects 5.8 million people, and each year 550,000 new cases are diagnosed. In 2011, heart failure was the most common reason for hospitalization for adults aged 85 years and older, and the second-most common for adults aged 65–84 years. An estimated one in five adults at age 40 will develop heart failure during their remaining lifetimes and about half of people who develop heart failure die within 5 years of diagnosis. Heart failure much higher in African Americans, Hispanics, Native Americans, and recent immigrants from Eastern Europe countries has been linked in these ethnic minority populations to the high incidence of diabetes and hypertension.
Nearly one of every four people (24.7%) hospitalized in the U.S. with congestive heart failure is readmitted within 30 days. Additionally, more than 50% of people seek readmission within 6 months after treatment and the average duration of hospital stay is 6 days. Heart failure is a leading cause of hospital readmissions in the U.S. People aged 65 and older were readmitted at a rate of 24.5 per 100 admissions in 2011. In the same year, heart failure patients under Medicaid were readmitted at a rate of 30.4 per 100 admissions, and uninsured people were readmitted at a rate of 16.8 per 100 admissions. These are the highest readmission rates for both categories. Notably, heart failure was not among the top-10 conditions with the most 30-day readmissions among the privately insured.
United Kingdom
In the UK, despite moderate improvements in prevention, heart failure rates have increased due to population growth and aging. Overall heart failure rates are similar to the four most common causes of cancer (breast, lung, prostate, and colon) combined. People from deprived backgrounds are more likely to be diagnosed with heart failure at a younger age.
Developing world
In tropical countries, the most common cause of heart failure is valvular heart disease or some type of cardiomyopathy. As underdeveloped countries have become more affluent, the incidences of diabetes, hypertension, and obesity have increased, which have in turn raised the incidence of heart failure.
Sex
Men have a higher incidence of heart failure, but the overall prevalence rate is similar in both sexes since women survive longer after the onset of heart failure. Women tend to be older when diagnosed with heart failure (after menopause), they are more likely than men to have diastolic dysfunction, and seem to experience a lower overall quality of life than men after diagnosis.
Ethnicity
Some sources state that people of Asian descent are at a higher risk of heart failure than other ethnic groups. Other sources however have found that rates of heart failure are similar to rates found in other ethnic groups.
History
For centuries, the disease entity which would include many cases of what today would be called heart failure was dropsy; the term denotes generalized edema, a major manifestation of a failing heart, though also caused by other diseases. Writings of ancient civilizations include evidence of their acquaintance with dropsy and heart failure: Egyptians were the first to use bloodletting to relieve fluid accumulation and shortage of breath, and provided what may have been the first documented observations on heart failure in the Ebers papyrus (around 1500 BCE). Greeks described cases of dyspnea, fluid retention and fatigue compatible with heart failure.
Romans used the flowering plant Drimia maritima (sea squill), which contains cardiac glycosides, for the treatment of dropsy; descriptions pertaining to heart failure are also known in the civilizations of ancient India and China. However, the manifestations of failing heart were understood in the context of these peoples' medical theories – including ancient Egyptian religion, Hippocratic theory of humours, or ancient Indian and Chinese medicine, and the current concept of heart failure had not developed yet. Although shortage of breath had been connected to heart disease by Avicenna round 1000 CE, decisive for modern understanding of the nature of the condition were the description of pulmonary circulation by Ibn al-Nafis in the 13th century, and of systemic circulation by William Harvey in 1628.
The role of the heart in fluid retention began to be better appreciated, as dropsy of the chest (fluid accumulation in and around the lungs causing shortness of breath) became more familiar and the current concept of heart failure, which brings together swelling and shortage of breath due to fluid retention, began to be accepted, in the 17th and especially in the 18th century: Richard Lower linked dyspnea and foot swelling in 1679, and Giovanni Maria Lancisi connected jugular vein distention with right ventricular failure in 1728. Dropsy attributable to other causes, e.g. kidney failure, was differentiated in the 19th century. The stethoscope, invented by René Laennec in 1819, x-rays, discovered by Wilhelm Röntgen in 1895, and electrocardiography, described by Willem Einthoven in 1903, facilitated the investigation of heart failure.
The 19th century also saw experimental and conceptual advances in the physiology of heart contraction, which led to the formulation of the Frank-Starling law of the heart (named after physiologists Otto Frank and Ernest Starling), a remarkable advance in understanding mechanisms of heart failure.
One of the earliest treatments of heart failure, relief of swelling by bloodletting with various methods, including leeches, continued through the centuries. Along with bloodletting, Jean-Baptiste de Sénac in 1749 recommended opiates for acute shortage of breath due to heart failure. In 1785, William Withering described the therapeutic uses of the foxglove genus of plants in the treatment of edema; their extract contains cardiac glycosides, including digoxin, still used today in the treatment of heart failure. The diuretic effects of inorganic mercury salts, which were used to treat syphilis, had already been noted in the 16th century by Paracelsus; in the 19th century they were used by noted physicians like John Blackall and William Stokes. In the meantime, cannulae (tubes) invented by English physician Reginald Southey in 1877 was another method of removing excess fluid by directly inserting into swollen limbs.
Use of organic mercury compounds as diuretics, beyond their role in syphilis treatment, started in 1920, though it was limited by their parenteral route of administration and their side-effects. Oral mercurial diuretics were introduced in the 1950s; so were thiazide diuretics, which caused less toxicity, and are still used. Around the same time, the invention of echocardiography by Inge Edler and Hellmuth Hertz in 1954 marked a new era in the evaluation of heart failure. In the 1960s, loop diuretics were added to available treatments of fluid retention, while a patient with heart failure received the first heart transplant by Christiaan Barnard. Over the following decades, new drug classes found their place in heart failure therapeutics, including vasodilators like hydralazine; renin-angiotensin system inhibitors; and beta-blockers.
Economics
In 2011, nonhypertensive heart failure was one of the 10 most expensive conditions seen during inpatient hospitalizations in the U.S., with aggregate inpatient hospital costs more than $10.5 billion.
Heart failure is associated with a high health expenditure, mostly because of the cost of hospitalizations; costs have been estimated to amount to 2% of the total budget of the National Health Service in the United Kingdom, and more than $35 billion in the United States.
Research directions
Some research indicates that stem cell therapy may help. Although this research indicated benefits of stem cell therapy, other research does not indicate benefit. There is tentative evidence of longer life expectancy and improved left ventricular ejection fraction in persons treated with bone marrow-derived stem cells.
The maintenance of heart function depends on appropriate gene expression that is regulated at multiple levels by epignetic mechanisms including DNA methylation and histone post-translational modification. Currently, an increasing body of research is directed at understanding the role of perturbations of epigenetic processes in cardiac hypertrophy and fibrotic scarring.
| Biology and health sciences | Symptoms and signs | Health |
249992 | https://en.wikipedia.org/wiki/Binary%20logarithm | Binary logarithm | In mathematics, the binary logarithm () is the power to which the number must be raised to obtain the value . That is, for any real number ,
For example, the binary logarithm of is , the binary logarithm of is , the binary logarithm of is , and the binary logarithm of is .
The binary logarithm is the logarithm to the base and is the inverse function of the power of two function. There are several alternatives to the notation for the binary logarithm; see the Notation section below.
Historically, the first application of binary logarithms was in music theory, by Leonhard Euler: the binary logarithm of a frequency ratio of two musical tones gives the number of octaves by which the tones differ. Binary logarithms can be used to calculate the length of the representation of a number in the binary numeral system, or the number of bits needed to encode a message in information theory. In computer science, they count the number of steps needed for binary search and related algorithms. Other areas
in which the binary logarithm is frequently used include combinatorics, bioinformatics, the design of sports tournaments, and photography.
Binary logarithms are included in the standard C mathematical functions and other mathematical software packages.
History
The powers of two have been known since antiquity; for instance, they appear in Euclid's Elements, Props. IX.32 (on the factorization of powers of two) and IX.36 (half of the Euclid–Euler theorem, on the structure of even perfect numbers).
And the binary logarithm of a power of two is just its position in the ordered sequence of powers of two.
On this basis, Michael Stifel has been credited with publishing the first known table of binary logarithms in 1544. His book Arithmetica Integra contains several tables that show the integers with their corresponding powers of two. Reversing the rows of these tables allow them to be interpreted as tables of binary logarithms.
Earlier than Stifel, the 8th century Jain mathematician Virasena is credited with a precursor to the binary logarithm. Virasena's concept of ardhacheda has been defined as the number of times a given number can be divided evenly by two. This definition gives rise to a function that coincides with the binary logarithm on the powers of two, but it is different for other integers, giving the 2-adic order rather than the logarithm.
The modern form of a binary logarithm, applying to any number (not just powers of two) was considered explicitly by Leonhard Euler in 1739. Euler established the application of binary logarithms to music theory, long before their applications in information theory and computer science became known. As part of his work in this area, Euler published a table of binary logarithms of the integers from 1 to 8, to seven decimal digits of accuracy.
Definition and properties
The binary logarithm function may be defined as the inverse function to the power of two function, which is a strictly increasing function over the positive real numbers and therefore has a unique inverse.
Alternatively, it may be defined as , where is the natural logarithm, defined in any of its standard ways. Using the complex logarithm in this definition allows the binary logarithm to be extended to the complex numbers.
As with other logarithms, the binary logarithm obeys the following equations, which can be used to simplify formulas that combine binary logarithms with multiplication or exponentiation:
For more, see list of logarithmic identities.
Notation
In mathematics, the binary logarithm of a number is often written as . However, several other notations for this function have been used or proposed, especially in application areas.
Some authors write the binary logarithm as , the notation listed in The Chicago Manual of Style. Donald Knuth credits this notation to a suggestion of Edward Reingold, but its use in both information theory and computer science dates to before Reingold was active. The binary logarithm has also been written as with a prior statement that the default base for the logarithm is . Another notation that is often used for the same function (especially in the German scientific literature) is , from Latin logarithmus dualis or logarithmus dyadis.
The , ISO 31-11 and ISO 80000-2 standards recommend yet another notation, . According to these standards, should not be used for the binary logarithm, as it is instead reserved for the common logarithm .
Applications
Information theory
The number of digits (bits) in the binary representation of a positive integer is the integral part of , i.e.
In information theory, the definition of the amount of self-information and information entropy is often expressed with the binary logarithm, corresponding to making the bit the fundamental unit of information. With these units, the Shannon–Hartley theorem expresses the information capacity of a channel as the binary logarithm of its signal-to-noise ratio, plus one. However, the natural logarithm and the nat are also used in alternative notations for these definitions.
Combinatorics
Although the natural logarithm is more important than the binary logarithm in many areas of pure mathematics such as number theory and mathematical analysis, the binary logarithm has several applications in combinatorics:
Every binary tree with leaves has height at least , with equality when is a power of two and the tree is a complete binary tree. Relatedly, the Strahler number of a river system with tributary streams is at most .
Every family of sets with different sets has at least elements in its union, with equality when the family is a power set.
Every partial cube with vertices has isometric dimension at least , and has at most edges, with equality when the partial cube is a hypercube graph.
According to Ramsey's theorem, every -vertex undirected graph has either a clique or an independent set of size logarithmic in . The precise size that can be guaranteed is not known, but the best bounds known on its size involve binary logarithms. In particular, all graphs have a clique or independent set of size at least and almost all graphs do not have a clique or independent set of size larger than .
From a mathematical analysis of the Gilbert–Shannon–Reeds model of random shuffles, one can show that the number of times one needs to shuffle an -card deck of cards, using riffle shuffles, to get a distribution on permutations that is close to uniformly random, is approximately . This calculation forms the basis for a recommendation that 52-card decks should be shuffled seven times.
Computational complexity
The binary logarithm also frequently appears in the analysis of algorithms, not only because of the frequent use of binary number arithmetic in algorithms, but also because binary logarithms occur in the analysis of algorithms based on two-way branching. If a problem initially has choices for its solution, and each iteration of the algorithm reduces the number of choices by a factor of two, then the number of iterations needed to select a single choice is again the integral part of . This idea is used in the analysis of several algorithms and data structures. For example, in binary search, the size of the problem to be solved is halved with each iteration, and therefore roughly iterations are needed to obtain a solution for a problem of size . Similarly, a perfectly balanced binary search tree containing elements has height .
The running time of an algorithm is usually expressed in big O notation, which is used to simplify expressions by omitting their constant factors and lower-order terms. Because logarithms in different bases differ from each other only by a constant factor, algorithms that run in time can also be said to run in, say, time. The base of the logarithm in expressions such as or is therefore not important and can be omitted. However, for logarithms that appear in the exponent of a time bound, the base of the logarithm cannot be omitted. For example, is not the same as because the former is equal to and the latter to .
Algorithms with running time are sometimes called linearithmic. Some examples of algorithms with running time or are:
Average time quicksort and other comparison sort algorithms
Searching in balanced binary search trees
Exponentiation by squaring
Longest increasing subsequence
Binary logarithms also occur in the exponents of the time bounds for some divide and conquer algorithms, such as the Karatsuba algorithm for multiplying -bit numbers in time ,
and the Strassen algorithm for multiplying matrices in time
. The occurrence of binary logarithms in these running times can be explained by reference to the master theorem for divide-and-conquer recurrences.
Bioinformatics
In bioinformatics, microarrays are used to measure how strongly different genes are expressed in a sample of biological material. Different rates of expression of a gene are often compared by using the binary logarithm of the ratio of expression rates: the log ratio of two expression rates is defined as the binary logarithm of the ratio of the two rates. Binary logarithms allow for a convenient comparison of expression rates: a doubled expression rate can be described by a log ratio of , a halved expression rate can be described by a log ratio of , and an unchanged expression rate can be described by a log ratio of zero, for instance.
Data points obtained in this way are often visualized as a scatterplot in which one or both of the coordinate axes are binary logarithms of intensity ratios, or in visualizations such as the MA plot and RA plot that rotate and scale these log ratio scatterplots.
Music theory
In music theory, the interval or perceptual difference between two tones is determined by the ratio of their frequencies. Intervals coming from rational number ratios with small numerators and denominators are perceived as particularly euphonious. The simplest and most important of these intervals is the octave, a frequency ratio of . The number of octaves by which two tones differ is the binary logarithm of their frequency ratio.
To study tuning systems and other aspects of music theory that require finer distinctions between tones, it is helpful to have a measure of the size of an interval that is finer than an octave and is additive (as logarithms are) rather than multiplicative (as frequency ratios are). That is, if tones , , and form a rising sequence of tones, then the measure of the interval from to plus the measure of the interval from to should equal the measure of the interval from to . Such a measure is given by the cent, which divides the octave into equal intervals ( semitones of cents each). Mathematically, given tones with frequencies and , the number of cents in the interval from to is
The millioctave is defined in the same way, but with a multiplier of instead of .
Sports scheduling
In competitive games and sports involving two players or teams in each game or match, the binary logarithm indicates the number of rounds necessary in a single-elimination tournament required to determine a winner. For example, a tournament of players requires rounds to determine the winner, a tournament of teams requires rounds, etc. In this case, for players/teams where is not a power of 2, is rounded up since it is necessary to have at least one round in which not all remaining competitors play. For example, is approximately , which rounds up to , indicating that a tournament of teams requires rounds (either two teams sit out the first round, or one team sits out the second round). The same number of rounds is also necessary to determine a clear winner in a Swiss-system tournament.
Photography
In photography, exposure values are measured in terms of the binary logarithm of the amount of light reaching the film or sensor, in accordance with the Weber–Fechner law describing a logarithmic response of the human visual system to light. A single stop of exposure is one unit on a base- logarithmic scale. More precisely, the exposure value of a photograph is defined as
where is the f-number measuring the aperture of the lens during the exposure, and is the number of seconds of exposure.
Binary logarithms (expressed as stops) are also used in densitometry, to express the dynamic range of light-sensitive materials or digital sensors.
Calculation
Conversion from other bases
An easy way to calculate on calculators that do not have a function is to use the natural logarithm () or the common logarithm ( or ) functions, which are found on most scientific calculators. To change the logarithm base to 2 from , , or any other base , one can use the formulae:
Approximately,
Integer rounding
The binary logarithm can be made into a function from integers and to integers by rounding it up or down. These two forms of integer binary logarithm are related by this formula:
The definition can be extended by defining . Extended in this way, this function is related to the number of leading zeros of the 32-bit unsigned binary representation of , .
The integer binary logarithm can be interpreted as the zero-based index of the most significant bit in the input. In this sense it is the complement of the find first set operation, which finds the index of the least significant bit. Many hardware platforms include support for finding the number of leading zeros, or equivalent operations, which can be used to quickly find the binary logarithm. The fls and flsl functions in the Linux kernel and in some versions of the libc software library also compute the binary logarithm (rounded up to an integer, plus one).
Iterative approximation
For a general positive real number, the binary logarithm may be computed in two parts.
First, one computes the integer part, (called the characteristic of the logarithm).
This reduces the problem to one where the argument of the logarithm is in a restricted range, the interval , simplifying the second step of computing the fractional part (the mantissa of the logarithm).
For any , there exists a unique integer such that , or equivalently . Now the integer part of the logarithm is simply , and the fractional part is . In other words:
For normalized floating-point numbers, the integer part is given by the floating-point exponent, and for integers it can be determined by performing a count leading zeros operation.
The fractional part of the result is and can be computed iteratively, using only elementary multiplication and division.
The algorithm for computing the fractional part can be described in pseudocode as follows:
Start with a real number in the half-open interval . If , then the algorithm is done, and the fractional part is zero.
Otherwise, square repeatedly until the result lies in the interval . Let be the number of squarings needed. That is, with chosen such that is in .
Taking the logarithm of both sides and doing some algebra:
Once again is a real number in the interval . Return to step 1 and compute the binary logarithm of using the same method.
The result of this is expressed by the following recursive formulas, in which is the number of squarings required in the i-th iteration of the algorithm:
In the special case where the fractional part in step 1 is found to be zero, this is a finite sequence terminating at some point. Otherwise, it is an infinite series that converges according to the ratio test, since each term is strictly less than the previous one (since every ). See Horner's method. For practical use, this infinite series must be truncated to reach an approximate result. If the series is truncated after the -th term, then the error in the result is less than .
Software library support
The log2 function is included in the standard C mathematical functions. The default version of this function takes double precision arguments but variants of it allow the argument to be single-precision or to be a long double. In computing environments supporting complex numbers and implicit type conversion such as MATLAB the argument to the log2 function is allowed to be a negative number, returning a complex one.
| Mathematics | Specific functions | null |
250046 | https://en.wikipedia.org/wiki/Common%20starling | Common starling | The common starling (Sturnus vulgaris), also known as simply as the starling in Great Britain and Ireland, and as European starling in North America, is a medium-sized passerine bird in the starling family, Sturnidae. It is about long and has glossy black plumage with a metallic sheen, which is speckled with white at some times of the year. The legs are pink and the bill is black in winter and yellow in summer; young birds have browner plumage than the adults. Its gift for mimicry has been noted in literature including the Mabinogion and the works of Pliny the Elder and William Shakespeare.
The common starling has about 12 subspecies breeding in open habitats across its native range in temperate Europe and across the Palearctic to western Mongolia, and it has been introduced as an invasive species to Australia, New Zealand, Canada, the United States, Mexico, Argentina, South Africa and Fiji. This bird is resident in western and southern Europe and southwestern Asia, while northeastern populations migrate south and west in the winter within the breeding range and also further south to Iberia and North Africa. The common starling builds an untidy nest in a natural or artificial cavity in which four or five glossy, pale blue eggs are laid. These take two weeks to hatch and the young remain in the nest for another three weeks. There are normally one or two breeding attempts each year. This species is omnivorous, taking a wide range of invertebrates, as well as seeds and fruit. It is hunted by various mammals and birds of prey, and is host to a range of external and internal parasites.
Large flocks typical of this species can be beneficial to agriculture by controlling invertebrate pests; however, starlings can also be pests themselves when they feed on fruit and sprouting crops. Common starlings may also be a nuisance through the noise and mess caused by their large urban roosts. Introduced populations in particular have been subjected to a range of controls, including culling, but these have had limited success, except in preventing the colonisation of Western Australia.
The species has declined in numbers in parts of northern and western Europe since the 1980s due to fewer grassland invertebrates being available as food for growing chicks. Despite this, its huge global population is not thought to be declining significantly, so the common starling is classified as being of least concern by the International Union for Conservation of Nature.
Taxonomy and systematics
The common starling was first described by Carl Linnaeus in his Systema Naturae in 1758 under its current binomial name. Sturnus and vulgaris are derived from the Latin for "starling" and "common" respectively. The Old English , later , and the Latin are both derived from an unknown Indo-European root dating back to the second millennium BC, imitative of the bird's call. "Starling" was first recorded in the 11th century, when it referred to the juvenile of the species, but by the 16th century it had already largely supplanted "stare" to refer to birds of all ages. The older name is referenced in William Butler Yeats' poem "The Stare's Nest by My Window". The International Ornithological Congress's preferred English vernacular name is common starling.
The starling family, Sturnidae, is an entirely Old World group apart from introductions elsewhere, with the greatest numbers of species in Southeast Asia and sub-Saharan Africa. The genus Sturnus is polyphyletic and relationships between its members are not fully resolved. The closest relation of the common starling is the spotless starling. The non-migratory spotless starling may be descended from a population of ancestral S. vulgaris that survived in an Iberian refugium during an Ice Age retreat, and mitochondrial gene studies suggest that it could be considered a subspecies of the common starling. There is more genetic variation between common starling populations than between the nominate common starling and the spotless starling. Although common starling remains are known from the Middle Pleistocene, part of the problem in resolving relationships in the Sturnidae is the paucity of the fossil record for the family as a whole.
Subspecies
There are several subspecies of the common starling, which vary clinally in size and the colour tone of the adult plumage. The gradual variation over geographic range and extensive intergradation means that acceptance of the various subspecies varies between authorities.
Birds from Fair Isle, St. Kilda and the Outer Hebrides are intermediate in size between S. v. zetlandicus and the nominate subspecies, and their subspecies placement varies according to the authority. The dark juveniles typical of these island forms are occasionally found in mainland Scotland and elsewhere, indicating some gene flow from S. v. faroensis or S. v. zetlandicus, subspecies formerly considered to be isolated.
Several other subspecies have been named, but are generally no longer considered valid. Most are intergrades that occur where the ranges of various subspecies meet. These include: S. v. ruthenus Menzbier, 1891 and S. v. jitkowi Buturlin, 1904, which are intergrades between S. v. vulgaris and S. v. poltaratskyi from western Russia; S. v. graecus Tschusi, 1905 and S. v. balcanicus Buturlin and Harms, 1909, which are intergrades between S. v. vulgaris and S. v. tauricus from the southern Balkans to central Ukraine and throughout Greece to the Bosporus; and S. v. heinrichi Stresemann, 1928, an intergrade between S. v. caucasicus and S. v. nobilior in northern Iran. S. v. persepolis Ticehurst, 1928 from southern Iran's (Fars province) is very similar to S. v. vulgaris, and it is not clear whether it is a distinct resident population or simply migrants from southeastern Europe.
Description
The common starling is long, with a wingspan of and a weight of . Among standard measurements, the wing chord is , the tail is , the culmen is and the tarsus is . The plumage is iridescent black, glossed purple or green, and spangled with white, especially in winter. The underparts of adult male common starlings are less spotted than those of adult females at a given time of year. The throat feathers of males are long and loose and are used in display while those of females are smaller and more pointed. The legs are stout, pinkish- or greyish-red in the breding season, and slightly darker in winter. The bill is narrow and conical with a sharp tip; in the winter it is brownish-black but in summer, females have lemon yellow beaks with pink bases while males have yellow bills with blue-grey bases. Moulting occurs once a year in late summer after the breeding season has finished; the fresh feathers are prominently tipped white (breast feathers) or buff (wing and back feathers), which gives the bird a speckled appearance. The reduction in the spotting in the breeding season is achieved through the white feather tips largely wearing off. Juveniles are grey-brown and by their first winter resemble adults though often retaining some brown juvenile feathering, especially on the head. They can usually be sexed by the colour of the irises, rich brown in males, mouse-brown or grey in females. Estimating the contrast between the iris and the central always-dark pupil is 97% accurate in determining sex, rising to 98% if the length of the throat feathers is also considered. The common starling is mid-sized by both starling standards and passerine standards. It is readily distinguished from most other mid-sized passerines, such as thrushes, icterids or small corvids, by its relatively short tail, sharp, blade-like bill, round-bellied shape and strong, sizeable (and rufous-coloured) legs. In flight, its strongly pointed wings and dark colour are distinctive, while on the ground its strange, somewhat waddling gait is also characteristic. The colouring and build usually distinguish this bird from other starlings, although the closely related spotless starling may be distinguished by the lack of pale feather tip spots in adult breeding plumage. The bohemian waxwing is structurally very similar in flight and also flies in dense flocks; it can be distinguished by being a paler reddish buff colour, marginally smaller, and also has a very different flight call.
Like most terrestrial starlings the common starling moves by walking or running, rather than hopping. Their flight is quite strong and direct; their triangular wings beat very rapidly, and periodically the birds glide for a short way without losing much height before resuming powered flight. When in a flock, the birds take off almost simultaneously, wheel and turn in unison, form a compact mass or trail off into a wispy stream, bunch up again and land in a coordinated fashion. Common starling on migration can fly at and cover up to .
Several terrestrial starlings, including those in the genus Sturnus, have adaptations of the skull and muscles that help with feeding by probing. This adaptation is most strongly developed in the common starling (along with the spotless and white-cheeked starlings), where the protractor muscles responsible for opening the jaw are enlarged and the skull is narrow, allowing the eye to be moved forward to peer down the length of the bill. This technique involves inserting the bill into the ground and opening it as a way of searching for hidden food items. Common starlings have the physical traits that enable them to use this feeding technique, which has undoubtedly helped the species spread far and wide.
In Iberia, the western Mediterranean and northwest Africa, the common starling may be confused with the closely related spotless starling, the plumage of which, as its name implies, has a more uniform colour. At close range it can be seen that the latter has longer throat feathers, a fact particularly noticeable when it sings.
Song and calls
The common starling is noisy, its song consisting of a wide variety of both melodic and mechanical-sounding noises as part of a ritual succession of sounds. The male is the main songster and engages in bouts of song lasting for a minute or more. Each of these typically includes four varieties of song type, which follow each other in a regular order without pause. The bout starts with a series of pure-tone whistles and these are followed by the main part of the song, a number of variable sequences that often incorporate snatches of song mimicked from other species of bird and various naturally occurring or man-made noises. The structure and simplicity of the sound mimicked is of greater importance than the frequency with which it occurs. In some instances, a wild starling has been observed to mimic a sound it has heard only once. Each sound clip is repeated several times before the bird moves on to the next. After this variable section comes a number of types of repeated clicks followed by a final burst of high-frequency song, again formed of several types. Each bird has its own repertoire with more proficient birds having a range of up to 35 variable song types and as many as 14 types of clicks.
Males sing constantly as the breeding period approaches and perform less often once pairs have bonded. In the presence of a female, a male sometimes flies to his nest and sings from the entrance, apparently attempting to entice the female in. Older birds tend to have a wider repertoire than younger ones. Those males that engage in longer bouts of singing and that have wider repertoires attract mates earlier and have greater reproductive success than others. Females appear to prefer mates with more complex songs, perhaps because this indicates greater experience or longevity. Having a complex song is also useful in defending a territory and deterring less experienced males from encroaching.Along with having adaptions of the skull and muscles for singing, male starlings also have a much larger syrinx than females. This is due to increased muscle mass and enlarged elements of the syringeal skeleton. The male starling's syrinx is around 35% larger than its female counterpart. However, this sexual dimorphism is less pronounced than it is in songbird species like the zebra finch, where the male's syrinx is 100% larger than the female's syrinx.
Singing also occurs outside the breeding season, taking place throughout the year apart from the moulting period. The songsters are more commonly male although females also sing on occasion. The function of such out-of-season song is poorly understood. Eleven other types of call have been described including a flock call, threat call, attack call, snarl call and copulation call. The alarm call is a harsh scream, and while foraging together common starlings squabble incessantly. They chatter while roosting and bathing, making a great deal of noise that can cause irritation to people living nearby. When a flock of common starlings is flying together, the synchronised movements of the birds' wings make a distinctive whooshing sound that can be heard hundreds of metres away.
Behaviour and ecology
The common starling is a highly gregarious species, especially in autumn and winter. Although flock size is highly variable, huge, noisy flocks (murmurations) may form near roosts. These dense concentrations of birds are thought to be a defence against attacks by birds of prey such as peregrine falcons or Eurasian sparrowhawks. Flocks form a tight sphere-like formation in flight, frequently expanding and contracting and changing shape, seemingly without any sort of leader. Each common starling changes its course and speed as a result of the movement of its closest neighbours.
Very large roosts, up to 1.5 million birds, form in city centres, woodlands and reedbeds, causing problems with their droppings. These may accumulate up to deep, killing trees by their concentration of chemicals. In smaller amounts, the droppings act as a fertiliser, and therefore woodland managers may try to move roosts from one area of a wood to another to benefit from the soil enhancement and avoid large toxic deposits.
Flocks of more than a million common starlings may be observed just before sunset in spring in southwestern Jutland, Denmark, over the seaward marshlands of Tønder and Esbjerg municipalities between Tønder and Ribe. They gather in March until northern Scandinavian birds leave for their breeding ranges by mid-April. Their swarm behaviour creates complex shapes silhouetted against the sky, a phenomenon known locally as sort sol ("black sun").
Flocks of anything from five to fifty thousand common starlings form in areas of the UK just before sundown during mid-winter. These flocks are commonly called murmurations.
Feeding
The common starling is largely insectivorous and feeds on both pest and other arthropods. The food range includes spiders, crane flies, moths, mayflies, dragonflies, damsel flies, grasshoppers, earwigs, lacewings, caddisflies, flies, beetles, sawflies, bees, wasps and ants. Prey are consumed in both adult and larvae stages of development, and common starlings will also feed on earthworms, snails, small amphibians and lizards. While the consumption of invertebrates is necessary for successful breeding, common starlings are omnivorous and can also eat grains, seeds, fruits, nectar and food waste if the opportunity arises. The Sturnidae differ from most birds in that they cannot easily metabolise foods containing high levels of sucrose, although they can cope with other fruits such as grapes and cherries. The isolated Azores subspecies of the common starling eats the eggs of the endangered roseate tern. Measures are being introduced to reduce common starling populations by culling before the terns return to their breeding colonies in spring.
There are several methods by which common starlings obtain their food, but, for the most part, they forage close to the ground, taking insects from the surface or just underneath. Generally, common starlings prefer foraging amongst short-cropped grasses and eat with grazing animals or perch on their backs, where they will also feed on the mammal's external parasites. Large flocks may engage in a practice known as "roller-feeding", where the birds at the back of the flock continually fly to the front where the feeding opportunities are best. The larger the flock, the nearer individuals are to one another while foraging. Flocks often feed in one place for some time, and return to previous successfully foraged sites.
There are three types of foraging behaviours observed in the common starling. "Probing" involves the bird plunging its beak into the ground randomly and repetitively until an insect has been found, and is often accompanied by bill gaping where the bird opens its beak in the soil to enlarge a hole. This behaviour, first described by Konrad Lorenz and given the German term zirkeln, is also used to create and widen holes in plastic garbage bags. It takes time for young common starlings to perfect this technique, and because of this the diet of young birds will often contain fewer insects. "Hawking" is the capture of flying insects directly from the air, and "lunging" is the less common technique of striking forward to catch a moving invertebrate on the ground. Earthworms are caught by pulling from soil. Common starlings that have periods without access to food, or have a reduction in the hours of light available for feeding, compensate by increasing their body mass by the deposition of fat.
Nesting
Unpaired males find a suitable cavity and begin to build nests in order to attract single females, often decorating the nest with ornaments such as flowers and fresh green material, which the female later disassembles upon accepting him as a mate. The amount of green material is not important, as long as some is present, but the presence of herbs in the decorative material appears to be significant in attracting a mate. The scent of plants such as yarrow acts as an olfactory attractant to females.
The males sing throughout much of the construction and even more so when a female approaches his nest. Following copulation, the male and female continue to build the nest. Nests may be in any type of hole, common locations include inside hollowed trees, buildings, tree stumps and man-made nest-boxes. S. v. zetlandicus typically breeds in crevices and holes in cliffs, a habitat only rarely used by the nominate form. Nests are typically made out of straw, dry grass and twigs with an inner lining made up of feathers, wool and soft leaves. Construction usually takes four or five days and may continue through incubation.
Common starlings are both monogamous and polygamous; although broods are generally brought up by one male and one female, occasionally the pair may have an extra helper. Pairs may be part of a colony, in which case several other nests may occupy the same or nearby trees. Males may mate with a second female while the first is still on the nest. The reproductive success of the bird is poorer in the second nest than it is in the primary nest and is better when the male remains monogamous.
Breeding
Breeding takes place during the spring and summer. Following copulation, the female lays eggs on a daily basis over a period of several days. If an egg is lost during this time, she will lay another to replace it. There are normally four or five eggs that are ovoid in shape and pale blue or occasionally white, and they commonly have a glossy appearance. The colour of the eggs seems to have evolved through the relatively good visibility of blue at low light levels. The egg size is in length and in maximum diameter.Incubation lasts thirteen days, although the last egg laid may take 24 hours longer than the first to hatch. Both parents share the responsibility of brooding the eggs, but the female spends more time incubating them than does the male, and is the only parent to do so at night when the male returns to the communal roost. A pair can raise up to three broods per year, frequently reusing and relining the same nest, although two broods is typical, or just one north of 48°N.
The young are born blind and naked. They develop light fluffy down within seven days of hatching and can see within nine days. As with other passerines, the nest is kept clean and the chicks' faecal sacs are removed by the adults. Once the chicks are able to regulate their body temperature, about six days after hatching, the adults largely cease removing droppings from the nest. Prior to that, the fouling would wet both the chicks' plumage and the nest material, thereby reducing their effectiveness as insulation and increasing the risk of chilling the hatchlings. Nestlings remain in the nest for three weeks, where they are fed continuously by both parents. Fledglings continue to be fed by their parents for another one or two weeks. Within two months, most juveniles will have moulted and gained their first basic plumage. They acquire their adult plumage the following year.
Intraspecific brood parasites are common in common starling nests. Female "floaters" (unpaired females during the breeding season) present in colonies often lay eggs in another pair's nest. Fledglings have also been reported to invade their own or neighbouring nests and evict a new brood. Common starling nests have a 48% to 79% rate of successful fledging, although only 20% of nestlings survive to breeding age; the adult survival rate is closer to 60%. The average life span is about 2–3 years, with a longevity record of 22 years 11 months.
Predators and parasites
A majority of starling predators are avian. The typical response of starling groups is to take flight, with a common sight being undulating flocks of starling flying high in quick and agile patterns. Their abilities in flight are seldom matched by birds of prey. Adult common starlings are hunted by hawks such as the northern goshawk (Accipiter gentilis) and Eurasian sparrowhawk (Accipiter nisus), and falcons including the peregrine falcon (Falco peregrinus), Eurasian hobby (Falco subbuteo) and common kestrel (Falco tinnunculus). Slower raptors like black and red kites (Milvus migrans & milvus), eastern imperial eagle (Aquila heliaca), common buzzard (Buteo buteo) and Australasian harrier (Circus approximans) tend to take the more easily caught fledglings or juveniles. While perched in groups by night, they can be vulnerable to owls, including the little owl (Athene noctua), long-eared owl (Asio otus), short-eared owl (Asio flammeus), barn owl (Tyto alba), tawny owl (Strix aluco) and Eurasian eagle-owl (Bubo bubo).
More than twenty species of hawk, owl and falcon are known to occasionally predate starlings in North America, though the most regular predators of adults are likely to be urban-living peregrine falcons or merlins (Falco columbarius). Common mynas (Acridotheres tristis) sometimes evict eggs, nestlings and adult common starlings from their nests, and the lesser honeyguide (Indicator minor), a brood parasite, uses the common starling as a host. Starlings are more commonly the culprits rather than victims of nest eviction however, especially towards other starlings and woodpeckers. Nests can be raided by mammals capable of climbing to them, such as small mustelids (Mustela spp.), raccoons (Procyon lotor) and squirrels (Sciurus spp.), and cats may catch the unwary.
Common starlings are hosts to a wide range of parasites. A survey of three hundred common starlings from six US states found that all had at least one type of parasite; 99% had external fleas, mites or ticks, and 95% carried internal parasites, mostly various types of worm. Blood-sucking species leave their host when it dies, but other external parasites stay on the corpse. A bird with a deformed bill was heavily infested with Mallophaga lice, presumably due to its inability to remove vermin.
The hen flea (Ceratophyllus gallinae) is the most common flea in their nests. The small, pale house-sparrow flea C. fringillae, is also occasionally found there and probably arises from the habit of its main host of taking over the nests of other species. This flea does not occur in the US, even on house sparrows. Lice include Menacanthus eurystemus, Brueelia nebulosa and Stumidoecus sturni. Other arthropod parasites include Ixodes ticks and mites such as Analgopsis passerinus, Boydaia stumi, Dermanyssus gallinae, Ornithonyssus bursa, O. sylviarum, Proctophyllodes species, Pteronyssoides truncatus and Trouessartia rosteri. The hen mite D. gallinae is itself preyed upon by the predatory mite Androlaelaps casalis. The presence of this control on numbers of the parasitic species may explain why birds are prepared to reuse old nests.
Flying insects that parasitise common starlings include the louse-fly Omithomya nigricornis and the saprophagous fly Carnus hemapterus. The latter species breaks off the feathers of its host and lives on the fats produced by growing plumage. Larvae of the moth Hofmannophila pseudospretella are nest scavengers, which feed on animal material such as faeces or dead nestlings. Protozoan blood parasites of the genus Haemoproteus have been found in common starlings, but a better known pest is the brilliant scarlet nematode Syngamus trachea. This worm moves from the lungs to the trachea and may cause its host to suffocate. In Britain, the rook and the common starling are the most infested wild birds. Other recorded internal parasites include the spiny-headed worm Prosthorhynchus transverses.
Common starlings may contract avian tuberculosis, avian malaria and retrovirus-induced lymphomas. Captive starlings often accumulate excess iron in the liver, a condition that can be prevented by adding black tea-leaves to the food.
Distribution and habitat
The global population of common starlings was estimated to be 310 million individuals in 2004, occupying a total area of . Widespread throughout the Northern Hemisphere, the bird is native to Eurasia and is found throughout Europe, northern Africa (from Morocco to Egypt), India (mainly in the north but regularly extending farther south and extending into the Maldives) Nepal, the Middle East including Israel, Syria, Iran, and Iraq, and northwestern China.
Common starlings in the south and west of Europe and south of latitude 40°N are mainly resident, although other populations migrate from regions where the winter is harsh, the ground frozen and food scarce. Large numbers of birds from northern Europe, Russia and Ukraine migrate south westwards or south eastwards. In the autumn, when immigrants are arriving from eastern Europe, many of Britain's common starlings are setting off for Iberia and North Africa. Other groups of birds are in passage across the country and the pathways of these different streams of bird may cross. Of the 15,000 birds ringed as nestlings in Merseyside, England, individuals have been recovered at various times of year as far afield as Norway, Sweden, Finland, Russia, Ukraine, Poland, Germany and the Low Countries. Small numbers of common starlings have sporadically been observed in Japan and Hong Kong but it is unclear whence these birds originated. In North America, northern populations have developed a migration pattern, vacating much of Canada in winter. Birds in the east of the country move southwards, and those from farther west winter in the southwest of the US.
Common starlings prefer urban or suburban areas where artificial structures and trees provide adequate nesting and roosting sites. Reedbeds are also favoured for roosting and the birds commonly feed in grassy areas such as farmland, grazing pastures, playing fields, golf courses and airfields where short grass makes foraging easy. They occasionally inhabit open forests and woodlands and are sometimes found in shrubby areas such as Australian heathland. Common starlings rarely inhabit dense, wet forests (i.e. rainforests or wet sclerophyll forests) but are found in coastal areas, where they nest and roost on cliffs and forage amongst seaweed. Their ability to adapt to a large variety of habitats has allowed them to disperse and establish themselves in diverse locations around the world resulting in a habitat range from coastal wetlands to alpine forests, from sea cliffs to mountain ranges above sea level.
Introduced populations
The common starling has been introduced to and has successfully established itself in New Zealand, Australia, South Africa, North America, Fiji and several Caribbean islands. As a result, it has also been able to migrate to Thailand, Southeast Asia and New Guinea.
South America
Five individuals conveyed on a ship from England alighted near Lago de Maracaibo in Venezuela in November 1949, but subsequently vanished. In 1987, a small population of common starlings was observed nesting in gardens in the city of Buenos Aires. Since then, despite some initial attempts at eradication, the bird has been expanding its breeding range at an average rate of per year, keeping within of the Atlantic coast. In Argentina, the species makes use of a variety of natural and man-made nesting sites, particularly woodpecker holes.
Australia
The common starling was introduced to Australia to consume insect pests of farm crops. Early settlers looked forward to their arrival, believing that common starlings were also important for the pollination of flax, a major agricultural product. Nest-boxes for the newly released birds were placed on farms and near crops. The common starling was introduced to Melbourne in 1857 and Sydney two decades later. By the 1880s, established populations were present in the southeast of the country thanks to the work of acclimatisation committees. By the 1920s, common starlings were widespread throughout Victoria, Queensland and New South Wales, but by then they were considered to be pests. Although common starlings were first sighted in Albany, Western Australia in 1917, they have been largely prevented from spreading to the state. The wide and arid Nullarbor Plain provides a natural barrier and control measures have been adopted that have killed 55,000 birds over three decades. The common starling has also colonised Kangaroo Island, Lord Howe Island, Norfolk Island and Tasmania.
New Zealand
The early settlers in New Zealand cleared the bush and found their newly planted crops were invaded by hordes of caterpillars and other insects deprived of their previous food sources. Native birds were not habituated to living in close proximity to man so the common starling was introduced from Europe along with the house sparrow to control the pests. It was first brought over in 1862 by the Nelson Acclimatisation Society and other introductions followed. The birds soon became established and are now found all over the country including the subtropical Kermadec Islands to the north and the equally distant Macquarie Island far to the south.
North America
Various acclimatisation society records mention instances of starlings being introduced in Cincinnati, Quebec and New York in the 1870s. As part of a nationwide effort, about 60 common starlings were released in 1890 into New York's Central Park by Eugene Schieffelin, president of the American Acclimatization Society. It has been widely reported that he had tried to introduce every bird species mentioned in the works of William Shakespeare into North America, but this claim has been traced to an essay in 1948 by naturalist Edwin Way Teale, whose notes appear to indicate that it was speculation. About the same date, the Portland Song Bird Club released 35 pairs of common starlings in Portland, Oregon. Earlier introductions are recorded to have died out within a few years, with the 1890 New York and Portland introductions reported as being the most successful.
Population of the birds is estimated to have grown to 150 million, occupying an area extending from southern Canada and Alaska to Central America.
Polynesia
The common starling appears to have arrived in Fiji in 1925 on Ono-i-lau and Vatoa islands. It may have colonised from New Zealand via Raoul in the Kermadec Islands where it is abundant, that group being roughly equidistant between New Zealand and Fiji. Its spread in Fiji has been limited, and there are doubts about the population's viability. Tonga was colonised at about the same date and the birds there have been slowly spreading north through the group.
South Africa
In South Africa, the common starling was introduced in 1897 by Cecil Rhodes. It spread slowly, and by 1954, had reached Clanwilliam and Port Elizabeth. It is now common in the southern Cape region, thinning out northwards to the Johannesburg area. It is present in the Western Cape, the Eastern Cape and the Free State provinces of South Africa and lowland Lesotho, with occasional sightings in KwaZulu-Natal, Gauteng and around the town of Oranjemund in Namibia. In Southern Africa populations appear to be resident and the bird is strongly associated with man and anthropogenic habitats. It favours irrigated land and is absent from regions where the ground is baked so dry that it cannot probe for insects. It may compete with native birds for crevice nesting sites, but the indigenous species are probably more disadvantaged by destruction of their natural habitat than they are by inter-specific competition. It breeds from September to December and outside the breeding season may congregate in large flocks, often roosting in reedbeds. It is the most common bird species in urban and agricultural areas.
West Indies
In 1901, the inhabitants of Saint Kitts petitioned the Colonial Secretary for a ″government grant of starlings to exterminate″ an outbreak of grasshoppers which was causing enormous damage to their crops. The common starling was introduced to Jamaica in 1903, and the Bahamas and Cuba were colonised naturally from the US. This bird is fairly common but local in Jamaica, Grand Bahama and Bimini, and is rare in the rest of the Bahamas and eastern Cuba.
Status
The global population of the common starling is estimated to be more than 310 million individuals and its numbers are not thought to be declining significantly, so the bird is classified by the International Union for Conservation of Nature as being of least concern. It had shown a marked increase in numbers throughout Europe from the 19th century to around the 1950s and 60s. In about 1830, S. v. vulgaris expanded its range in the British Isles, spreading into Ireland and areas of Scotland where it had formerly been absent, although S. v. zetlandicus was already present in Shetland and the Outer Hebrides. The common starling has bred in northern Sweden from 1850 and in Iceland from 1935. The breeding range spread through southern France to northeastern Spain, and there were other range expansions particularly in Italy, Austria and Finland. It started breeding in Iberia in 1960, while the spotless starling's range had been expanding northward since the 1950s. The low rate of advance, about per year for both species, is due to the suboptimal mountain and woodland terrain. Expansion has since slowed even further due to direct competition between the two similar species where they overlap in southwestern France and northwestern Spain.
Major declines in populations have been observed from 1980 onward in Sweden, Finland, northern Russia (Karelia) and the Baltic States, and smaller declines in much of the rest of northern and central Europe. The bird has been adversely affected in these areas by intensive agriculture, and in several countries it has been red-listed due to population declines of more than 50%. Numbers dwindled in the United Kingdom by more than 80% between 1966 and 2004; although populations in some areas such as Northern Ireland were stable or even increased, those in other areas, mainly England, declined even more sharply. The overall decline seems to be due to the low survival rate of young birds, which may be caused by changes in agricultural practices. The intensive farming methods used in northern Europe mean there is less pasture and meadow habitat available, and the supply of grassland invertebrates needed for the nestlings to thrive is correspondingly reduced.
Relationship with humans
Benefits and problems
Since common starlings eat insect pests such as wireworms, they are considered beneficial in northern Eurasia, and this was one of the reasons given for introducing the birds elsewhere. Around 25 million nest boxes were erected for this species in the former Soviet Union, and common starlings were found to be effective in controlling the grass grub Costelytra zealandica in New Zealand. The original Australian introduction was facilitated by the provision of nest boxes to help this mainly insectivorous bird to breed successfully, and even in the US, where this is a pest species, the Department of Agriculture acknowledges that vast numbers of insects are consumed by common starlings.
Common starlings introduced to areas such as Australia or North America, where other members of the genus are absent, may affect native species through competition for nest holes. In North America, chickadees, nuthatches, woodpeckers, purple martins and other swallows may be affected. In Australia, competitors for nesting sites include the crimson and eastern rosellas. For its role in the decline of local native species and the damages to agriculture, the common starling has been included in the IUCN List of the world's 100 worst invasive species.European, or common, starlings are habitat generalists meaning they are able to exploit a multitude of habitats, nest sites and food sources. This, coupled with them being lowland birds that easily coexist with humans, enables them to take advantage of other native birds, most particularly woodpecker. European starlings are considered aggressive omnivores that utilize an open-bill probing technique that gives them an evolutionary advantage over birds that are frugivores. Their aggressive and gregarious behaviour in terms of food thus allows them to outcompete native species. Common starlings are also aggressive in the creation of their nest cavities. Often, starlings will usurp a nest site, for example a tree hollow, and fill it rapidly with bedding and contaminants compared to other species, like the native parrots, that use little to no bedding. As cavity nesters, they are able to outcompete many native species in terms of habitat and nest sites.
Common starlings can eat and damage fruit in orchards such as grapes, peaches, olives, currants and tomatoes or dig up newly sown grain and sprouting crops. They may also eat animal feed and distribute seeds through their droppings. In eastern Australia, weeds like bridal creeper, blackberry and boneseed are thought to have been spread by common starlings. Agricultural damage in the US is estimated as costing about US$800million annually. This bird is not considered to be as damaging to agriculture in South Africa as it is in the United States.
Common starlings take advantage of agricultural fields, livestock facilities, and other human related sources of food and nest sites. Starlings often assault crops such as grapes, olives, and cherries by consuming excessive amounts of crops in large flock sizes and in new grain fields, starlings pull up young plants and eat the seeds. In caged trials, it was shown that starlings eat of animal food daily and of plant food meaning a decent portion of crops are consumed by these birds. Bird damage to grapes in 1968 cost upwards to $4.4million while losing almost 17% of the crops. Common starlings also often congregate at feeding troughs to eat grain and concurrently contaminate the food and water sources provided for livestock with their droppings. For example, high protein supplements added to cattle feed are selectively eaten by common starlings. In 1968, the cost of cattle rations consumed during winter by starlings was $84 per 1,000 starlings and is proposed to be much more expensive today given an increase in current cattle feed costs. The English or house sparrow (Passer domesticus) and the common starling are considerable agricultural pests, together causing an estimated US$1billion per year in crop damages.
The large size of flocks can also cause problems. The large roosts of the common starling pose many safety hazards for aircraft, mainly including the clogging of engines that concurrently shutdown the plane into descent. One of the worst instances of this was the Eastern Air Lines Flight 375 incident in Boston in 1960, when 62 people died after a turboprop airliner flew into a flock and plummeted into the sea at Winthrop Harbor. From the years 1990–2001, 852 incidents of aircraft hazard due to starlings and New World blackbirds were reported with 39 strikes causing major damage that cost a total of $1,607,317.
Starlings' droppings can contain the fungus Histoplasma capsulatum, the cause of histoplasmosis in humans. At roosting sites this fungus can thrive in accumulated droppings. There are a number of other infectious diseases that can potentially be transmitted by common starlings to humans, although the potential for the birds to spread infections may have been exaggerated. The spread of disease to livestock is also a concern, possibly more important than starling's effects on food consumption or transmission of disease to humans. The spreading of Histoplasmosis reported in a Nebraska manufacturing facility saw a loss of 10,000 pigs from the spread of the disease which was valued at $1million loss in 2014.
Control
Due to the impact of starlings on crop production, there have been attempts to control the numbers of both native and introduced populations of common starlings. Within the natural breeding range, this may be affected by legislation. For example, in Spain, the species is hunted commercially as a food item, and has a closed season, whereas in France, it is classed as a pest, and the season in which it may be killed covers the greater part of the year. In Great Britain, starlings are protected under the Wildlife and Countryside Act 1981, which makes it "illegal to intentionally kill, injure or take a starling, or to take, damage or destroy an active nest or its contents". The Wildlife Order in Northern Ireland allows, with a general licence, "an authorised person to control starlings to prevent serious damage to agriculture or preserve public health and safety". The species is migratory, so birds involved in control measures may have come from a wide area and breeding populations may not be greatly affected. In Europe, the varying legislation and mobile populations mean that control attempts may have limited long-term results. Non-lethal techniques such as scaring with visual or auditory devices have only a temporary effect in any case.
Huge urban roosts in cities can create problems due to the noise and mess made and the smell of the droppings. In 1949, so many birds landed on the clock hands of London's Big Ben that it stopped, leading to unsuccessful attempts to disrupt the roosts with netting, repellent chemical on the ledges and broadcasts of common starling alarm calls. An entire episode of The Goon Show in 1954 was a parody of the futile efforts to disrupt the large common starling roosts in central London.
Where it is introduced, the common starling is unprotected by legislation, and extensive control plans may be initiated. Common starlings can be prevented from using nest boxes by ensuring that the access holes are smaller than the diameter they need, and the removal of perches discourages them from visiting bird feeders.
Western Australia banned the import of common starlings in 1895. New flocks arriving from the east are routinely shot, while the less cautious juveniles are trapped and netted. New methods are being developed, such as tagging one bird and tracking it back to establish where other members of the flock roost. Another technique is to analyse the DNA of Australian common starling populations to track where the migration from eastern to western Australia is occurring so that better preventive strategies can be used. By 2009, only 300 common starlings were left in Western Australia, and the state committed a further A$400,000 in that year to continue the eradication programme.
In the United States, common starlings are exempt from the Migratory Bird Treaty Act, which prohibits the taking or killing of migratory birds. No permit is required to remove nests and eggs or kill juveniles or adults. Research was undertaken in 1966 to identify a suitable avicide that would both kill common starlings and would readily be eaten by them. It also needed to be of low toxicity to mammals and not likely to cause the death of pets that ate dead birds. The chemical that best fitted these criteria was DRC-1339, now marketed as Starlicide. In 2008, the United States government poisoned, shot or trapped 1.7million birds, the largest number of any nuisance species to be culled. In 2005, the population in the United States was estimated at 140million birds, around 45% of the global total of 310million.
The likelihood of starlings to damage the feeding operations is dependent on the number of livestock, favouring areas with more livestock. They also show preference for feed types which were not whole corn but smaller feeds, creating more damage in areas where the feed was smaller. They also showed feed preference based on composition. A proposed solution to this problem is use of less palatable feed by agriculturalists, perhaps relying on larger feed types or feed which is less favourable in composition to starlings. An additional solution for mitigation control involves ensuring that livestock feeding operations are not within close proximity of each other or starling roosts. Weather conditions also had an impact on whether starlings visited livestock feeding operations, with a higher likelihood to visit in colder temperatures or following snow storms.
Alternatives to managing starling populations in agricultural areas include the use of starlicide. Use of starlicide has been found to reduce the spread of Salmonella enterica in livestock and other diseases found among livestock. Though this does not appear to eliminate introduction of these diseases completely, it has been determined that they are contributors and starling control is a successful mitigation strategy.
In science and culture
Common starlings may be kept as pets or as laboratory animals. Austrian ethologist Konrad Lorenz wrote of them in his book King Solomon's Ring as "the poor man's dog" and "something to love", because nestlings are easily obtained from the wild and after careful hand rearing they are straightforward to look after. They adapt well to captivity, and thrive on a diet of standard bird feed and mealworms. Several birds may be kept in the same cage, and their inquisitiveness makes them easy to train or study. The only disadvantages are their messy and indiscriminate defecation habits and the need to take precautions against diseases that may be transmitted to humans. As a laboratory bird, the common starling is second in numbers only to the domestic pigeon.
The common starling's gift for mimicry has long been recognised. In the medieval Welsh , Branwen tamed a common starling, "taught it words", and sent it across the Irish Sea with a message to her brothers, Bran and Manawydan, who then sailed from Wales to Ireland to rescue her. Pliny the Elder claimed that these birds could be taught to speak whole sentences in Latin and Greek, and in Henry IV, William Shakespeare had Hotspur declare "The king forbade my tongue to speak of Mortimer. But I will find him when he is asleep, and in his ear I'll holler 'Mortimer!' Nay I'll have a starling shall be taught to speak nothing but Mortimer, and give it to him to keep his anger still in motion."
Mozart had a pet common starling which could sing part of his Piano Concerto in G Major (KV. 453). He had bought it from a shop after hearing it sing a phrase from a work he wrote six weeks previously, which had not yet been performed in public. He became very attached to the bird and arranged an elaborate funeral for it when it died three years later. It has been suggested that his A Musical Joke (K. 522) might be written in the comical, inconsequential style of a starling's vocalisation. Other people who have owned common starlings report how adept they are at picking up phrases and expressions. The words have no meaning for the starling, so they often mix them up or use them on what to humans are inappropriate occasions in their songs. Their ability at mimicry is so great that strangers have looked in vain for the human they think they have just heard speak.
Common starlings are trapped for food in some Arab countries. The meat is tough and of low quality, so it is casseroled or made into pâté. One recipe said it should be stewed "until tender, however long that may be". Even when correctly prepared, it may still be seen as an acquired taste.
| Biology and health sciences | Passerida | null |
250424 | https://en.wikipedia.org/wiki/Ring%20theory | Ring theory | In algebra, ring theory is the study of rings, algebraic structures in which addition and multiplication are defined and have similar properties to those operations defined for the integers. Ring theory studies the structure of rings; their representations, or, in different language, modules; special classes of rings (group rings, division rings, universal enveloping algebras); related structures like rngs; as well as an array of properties that prove to be of interest both within the theory itself and for its applications, such as homological properties and polynomial identities.
Commutative rings are much better understood than noncommutative ones. Algebraic geometry and algebraic number theory, which provide many natural examples of commutative rings, have driven much of the development of commutative ring theory, which is now, under the name of commutative algebra, a major area of modern mathematics. Because these three fields (algebraic geometry, algebraic number theory and commutative algebra) are so intimately connected it is usually difficult and meaningless to decide which field a particular result belongs to. For example, Hilbert's Nullstellensatz is a theorem which is fundamental for algebraic geometry, and is stated and proved in terms of commutative algebra. Similarly, Fermat's Last Theorem is stated in terms of elementary arithmetic, which is a part of commutative algebra, but its proof involves deep results of both algebraic number theory and algebraic geometry.
Noncommutative rings are quite different in flavour, since more unusual behavior can arise. While the theory has developed in its own right, a fairly recent trend has sought to parallel the commutative development by building the theory of certain classes of noncommutative rings in a geometric fashion as if they were rings of functions on (non-existent) 'noncommutative spaces'. This trend started in the 1980s with the development of noncommutative geometry and with the discovery of quantum groups. It has led to a better understanding of noncommutative rings, especially noncommutative Noetherian rings.
For the definitions of a ring and basic concepts and their properties, see Ring (mathematics). The definitions of terms used throughout ring theory may be found in Glossary of ring theory.
Commutative rings
A ring is called commutative if its multiplication is commutative. Commutative rings resemble familiar number systems, and various definitions for commutative rings are designed to formalize properties of the integers. Commutative rings are also important in algebraic geometry. In commutative ring theory, numbers are often replaced by ideals, and the definition of the prime ideal tries to capture the essence of prime numbers. Integral domains, non-trivial commutative rings where no two non-zero elements multiply to give zero, generalize another property of the integers and serve as the proper realm to study divisibility. Principal ideal domains are integral domains in which every ideal can be generated by a single element, another property shared by the integers. Euclidean domains are integral domains in which the Euclidean algorithm can be carried out. Important examples of commutative rings can be constructed as rings of polynomials and their factor rings. Summary: Euclidean domain ⊂ principal ideal domain ⊂ unique factorization domain ⊂ integral domain ⊂ commutative ring.
Algebraic geometry
Algebraic geometry is in many ways the mirror image of commutative algebra. This correspondence started with Hilbert's Nullstellensatz that establishes a one-to-one correspondence between the points of an algebraic variety, and the maximal ideals of its coordinate ring. This correspondence has been enlarged and systematized for translating (and proving) most geometrical properties of algebraic varieties into algebraic properties of associated commutative rings. Alexander Grothendieck completed this by introducing schemes, a generalization of algebraic varieties, which may be built from any commutative ring. More precisely,
the spectrum of a commutative ring is the space of its prime ideals equipped with Zariski topology, and augmented with a sheaf of rings. These objects are the "affine schemes" (generalization of affine varieties), and a general scheme is then obtained by "gluing together" (by purely algebraic methods) several such affine schemes, in analogy to the way of constructing a manifold by gluing together the charts of an atlas.
Noncommutative rings
Noncommutative rings resemble rings of matrices in many respects. Following the model of algebraic geometry, attempts have been made recently at defining noncommutative geometry based on noncommutative rings.
Noncommutative rings and associative algebras (rings that are also vector spaces) are often studied via their categories of modules. A module over a ring is an abelian group that the ring acts on as a ring of endomorphisms, very much akin to the way fields (integral domains in which every non-zero element is invertible) act on vector spaces. Examples of noncommutative rings are given by rings of square matrices or more generally by rings of endomorphisms of abelian groups or modules, and by monoid rings.
Representation theory
Representation theory is a branch of mathematics that draws heavily on non-commutative rings. It studies abstract algebraic structures by representing their elements as linear transformations of vector spaces, and studies
modules over these abstract algebraic structures. In essence, a representation makes an abstract algebraic object more concrete by describing its elements by matrices and the algebraic operations in terms of matrix addition and matrix multiplication, which is non-commutative. The algebraic objects amenable to such a description include groups, associative algebras and Lie algebras. The most prominent of these (and historically the first) is the representation theory of groups, in which elements of a group are represented by invertible matrices in such a way that the group operation is matrix multiplication.
Some relevant theorems
General
Isomorphism theorems for rings
Nakayama's lemma
Structure theorems
The Artin–Wedderburn theorem determines the structure of semisimple rings
The Jacobson density theorem determines the structure of primitive rings
Goldie's theorem determines the structure of semiprime Goldie rings
The Zariski–Samuel theorem determines the structure of a commutative principal ideal ring
The Hopkins–Levitzki theorem gives necessary and sufficient conditions for a Noetherian ring to be an Artinian ring
Morita theory consists of theorems determining when two rings have "equivalent" module categories
Cartan–Brauer–Hua theorem gives insight on the structure of division rings
Wedderburn's little theorem states that finite domains are fields
Other
The Skolem–Noether theorem characterizes the automorphisms of simple rings
Structures and invariants of rings
Dimension of a commutative ring
In this section, R denotes a commutative ring. The Krull dimension of R is the supremum of the lengths n of all the chains of prime ideals . It turns out that the polynomial ring over a field k has dimension n. The fundamental theorem of dimension theory states that the following numbers coincide for a noetherian local ring :
The Krull dimension of R.
The minimum number of the generators of the -primary ideals.
The dimension of the graded ring (equivalently, 1 plus the degree of its Hilbert polynomial).
A commutative ring R is said to be catenary if for every pair of prime ideals , there exists a finite chain of prime ideals that is maximal in the sense that it is impossible to insert an additional prime ideal between two ideals in the chain, and all such maximal chains between and have the same length. Practically all noetherian rings that appear in applications are catenary. Ratliff proved that a noetherian local integral domain R is catenary if and only if for every prime ideal ,
where is the height of .
If R is an integral domain that is a finitely generated k-algebra, then its dimension is the transcendence degree of its field of fractions over k. If S is an integral extension of a commutative ring R, then S and R have the same dimension.
Closely related concepts are those of depth and global dimension. In general, if R is a noetherian local ring, then the depth of R is less than or equal to the dimension of R. When the equality holds, R is called a Cohen–Macaulay ring. A regular local ring is an example of a Cohen–Macaulay ring. It is a theorem of Serre that R is a regular local ring if and only if it has finite global dimension and in that case the global dimension is the Krull dimension of R. The significance of this is that a global dimension is a homological notion.
Morita equivalence
Two rings R, S are said to be Morita equivalent if the category of left modules over R is equivalent to the category of left modules over S. In fact, two commutative rings which are Morita equivalent must be isomorphic, so the notion does not add anything new to the category of commutative rings. However, commutative rings can be Morita equivalent to noncommutative rings, so Morita equivalence is coarser than isomorphism. Morita equivalence is especially important in algebraic topology and functional analysis.
Finitely generated projective module over a ring and Picard group
Let R be a commutative ring and the set of isomorphism classes of finitely generated projective modules over R; let also subsets consisting of those with constant rank n. (The rank of a module M is the continuous function .) is usually denoted by Pic(R). It is an abelian group called the Picard group of R. If R is an integral domain with the field of fractions F of R, then there is an exact sequence of groups:
where is the set of fractional ideals of R. If R is a regular domain (i.e., regular at any prime ideal), then Pic(R) is precisely the divisor class group of R.
For example, if R is a principal ideal domain, then Pic(R) vanishes. In algebraic number theory, R will be taken to be the ring of integers, which is Dedekind and thus regular. It follows that Pic(R) is a finite group (finiteness of class number) that measures the deviation of the ring of integers from being a PID.
One can also consider the group completion of ; this results in a commutative ring K0(R). Note that K0(R) = K0(S) if two commutative rings R, S are Morita equivalent.
Structure of noncommutative rings
The structure of a noncommutative ring is more complicated than that of a commutative ring. For example, there exist simple rings that contain no non-trivial proper (two-sided) ideals, yet contain non-trivial proper left or right ideals. Various invariants exist for commutative rings, whereas invariants of noncommutative rings are difficult to find. As an example, the nilradical of a ring, the set of all nilpotent elements, is not necessarily an ideal unless the ring is commutative. Specifically, the set of all nilpotent elements in the ring of all matrices over a division ring never forms an ideal, irrespective of the division ring chosen. There are, however, analogues of the nilradical defined for noncommutative rings, that coincide with the nilradical when commutativity is assumed.
The concept of the Jacobson radical of a ring; that is, the intersection of all right (left) annihilators of simple right (left) modules over a ring, is one example. The fact that the Jacobson radical can be viewed as the intersection of all maximal right (left) ideals in the ring, shows how the internal structure of the ring is reflected by its modules. It is also a fact that the intersection of all maximal right ideals in a ring is the same as the intersection of all maximal left ideals in the ring, in the context of all rings; irrespective of whether the ring is commutative.
Noncommutative rings are an active area of research due to their ubiquity in mathematics. For instance, the ring of n-by-n matrices over a field is noncommutative despite its natural occurrence in geometry, physics and many parts of mathematics. More generally, endomorphism rings of abelian groups are rarely commutative, the simplest example being the endomorphism ring of the Klein four-group.
One of the best-known strictly noncommutative ring is the quaternions.
Applications
The ring of integers of a number field
The coordinate ring of an algebraic variety
If X is an affine algebraic variety, then the set of all regular functions on X forms a ring called the coordinate ring of X. For a projective variety, there is an analogous ring called the homogeneous coordinate ring. Those rings are essentially the same things as varieties: they correspond in essentially a unique way. This may be seen via either Hilbert's Nullstellensatz or scheme-theoretic constructions (i.e., Spec and Proj).
Ring of invariants
A basic (and perhaps the most fundamental) question in the classical invariant theory is to find and study polynomials in the polynomial ring that are invariant under the action of a finite group (or more generally reductive) G on V. The main example is the ring of symmetric polynomials: symmetric polynomials are polynomials that are invariant under permutation of variable. The fundamental theorem of symmetric polynomials states that this ring is where are elementary symmetric polynomials.
History
Commutative ring theory originated in algebraic number theory, algebraic geometry, and invariant theory. Central to the development of these subjects were the rings of integers in algebraic number fields and algebraic function fields, and the rings of polynomials in two or more variables. Noncommutative ring theory began with attempts to extend the complex numbers to various hypercomplex number systems. The genesis of the theories of commutative and noncommutative rings dates back to the early 19th century, while their maturity was achieved only in the third decade of the 20th century.
More precisely, William Rowan Hamilton put forth the quaternions and biquaternions; James Cockle presented tessarines and coquaternions; and William Kingdon Clifford was an enthusiast of split-biquaternions, which he called algebraic motors. These noncommutative algebras, and the non-associative Lie algebras, were studied within universal algebra before the subject was divided into particular mathematical structure types. One sign of re-organization was the use of direct sums to describe algebraic structure.
The various hypercomplex numbers were identified with matrix rings by Joseph Wedderburn (1908) and Emil Artin (1928). Wedderburn's structure theorems were formulated for finite-dimensional algebras over a field while Artin generalized them to Artinian rings.
In 1920, Emmy Noether, in collaboration with W. Schmeidler, published a paper about the theory of ideals in which they defined left and right ideals in a ring. The following year she published a landmark paper called Idealtheorie in Ringbereichen, analyzing ascending chain conditions with regard to (mathematical) ideals. Noted algebraist Irving Kaplansky called this work "revolutionary"; the publication gave rise to the term "Noetherian ring", and several other mathematical objects being called Noetherian.
| Mathematics | Algebra | null |
250460 | https://en.wikipedia.org/wiki/Potassium%20bromide | Potassium bromide | Potassium bromide (KBr) is a salt, widely used as an anticonvulsant and a sedative in the late 19th and early 20th centuries, with over-the-counter use extending to 1975 in the US. Its action is due to the bromide ion (sodium bromide is equally effective). Potassium bromide is used as a veterinary drug, in antiepileptic medication for dogs.
Under standard conditions, potassium bromide is a white crystalline powder. It is freely soluble in water; it is not soluble in acetonitrile. In a dilute aqueous solution, potassium bromide tastes sweet, at higher concentrations it tastes bitter, and tastes salty when the concentration is even higher. These effects are mainly due to the properties of the potassium ion—sodium bromide tastes salty at any concentration. In high concentration, potassium bromide strongly irritates the gastric mucous membrane, causing nausea and sometimes vomiting (a typical effect of all soluble potassium salts).
Chemical properties
Potassium bromide, a typical ionic salt, is fully dissociated and near pH 7 in aqueous solution. It serves as a source of bromide ions. This reaction is important for the manufacture of silver bromide for photographic film:
KBr_{(aq)}{} + AgNO3_{(aq)} -> AgBr_{(s)}{} + KNO3_{(aq)}
Aqueous bromide also forms complexes when reacted with some metal halides such as copper(II) bromide:
2 KBr_{(aq)}{} + CuBr2_{(aq)} -> K2[CuBr4]_{(aq)}
Preparation
A traditional method for the manufacture of KBr is the reaction of potassium carbonate with an iron(III, II) bromide, , made by treating scrap iron under water with excess bromine:
4 K2CO3 + Fe3Br8 -> 8 KBr + Fe3O4 + 4 CO2
Applications
Medical and veterinary
The anticonvulsant properties of potassium bromide were first noted by Sir Charles Locock at a meeting of the Royal Medical and Chirurgical Society in 1857. Bromide can be regarded as the first effective medication for epilepsy. At the time, it was commonly thought that epilepsy was caused by masturbation. Locock noted that bromide calmed sexual excitement and thought this was responsible for his success in treating seizures. In the latter half of the 19th century, potassium bromide was used for the calming of seizure and nervous disorders on an enormous scale, with the use by single hospitals being as much as several tons a year (the dose for a given person being a few grams per day). By the beginning of the 20th century the generic word had become so widely associated with being sedate that bromide came to mean a dull, sedate person or a boring platitude uttered by such a person.
There was not a better epilepsy drug until phenobarbital in 1912. The British Army has historically been claimed to lace soldiers' tea with bromide to quell sexual arousal and in the Victorian era prisoners in England were compulsorily dosed with the chemical.
Bromide compounds, especially sodium bromide, remained in over-the-counter sedatives and headache remedies (such as the original formulation of Bromo-Seltzer) in the US until 1975, when bromides were outlawed in all over-the-counter medicines, due to chronic toxicity. Bromide's exceedingly long half life in the body made it difficult to dose without side effects. Medical use of bromides in the US was discontinued at this time, as many better and shorter-acting sedatives were known by then.
Potassium bromide is used in veterinary medicine to treat epilepsy in dogs, either as first-line treatment or in addition to phenobarbital, when seizures are not adequately controlled with phenobarbital alone. Use of bromide in cats is limited because it carries a substantial risk of causing lung inflammation (pneumonitis) in them. Why bromides should cause such inflammation in cats, but not in dogs is not clear.
The use of bromide as a treatment drug for animals means that veterinary medical diagnostic laboratories are able as a matter of routine to measure serum levels of bromide on order of a veterinarian, whereas human medical diagnostic labs in the US do not measure bromide as a routine test.
Potassium bromide is not approved by the US Food and Drug Administration (FDA) for use in humans to control seizures. In Germany, it is still approved as an antiepileptic drug for humans, particularly children and adolescents. These indications include severe forms of generalized tonic-clonic seizures, early-childhood-related tonic–clonic seizures, and also severe myoclonic seizures during childhood. Adults who have reacted positively to the drug during childhood/adolescence may continue treatment. Potassium bromide tablets are sold under the brand name Dibro-Be mono (Rx-only). The drug has almost complete bioavailability, but the bromide ion has a relatively long half life of 12 days in the blood, making bromide salts difficult to adjust and dose. Bromide is not known to interfere with the absorption or excretion of any other anticonvulsant, though it does have strong interactions with chloride in the body, the normal body uptake and excretion of which strongly influences bromide's excretion.
The therapeutic index (ratio of effectiveness to toxicity) for bromide is small. As with other antiepileptics, sometimes even therapeutic doses (3 to 5 grams per day, taking 6 to 8 weeks to reach stable levels) may give rise to intoxication. Often indistinguishable from 'expected' side-effects, these include:
Bromism These are central nervous system reactions. They may include:
depression,
lethargy, somnolence (from daytime sleepiness to coma)
loss of appetite and cachexia, nausea/emesis with exicosis (loss of body fluid)
loss of reflexes or pathologic reflexes
clonic seizures
tremor
ataxia
loss of neural sensitivity
paresis
cerebral edema with associated headache and papilledema of the eyes
delirium: confusion, abnormal speech, loss of concentration and memory, aggressiveness
psychosis
Acne-form dermatitis and other forms of skin disease may also be seen, as well as mucous hypersecretion in the lungs. Asthma and rhinitis may worsen. Rarely, tongue disorder, aphthous stomatitis, bad breath, and constipation occur.
Optics
Potassium bromide is transparent from the near ultraviolet to long-wave infrared wavelengths (0.25-25 μm) and has no significant optical absorption lines in its high transmission region.
It is used widely as infrared optical windows and components for general spectroscopy because of its wide spectral range. In infrared spectroscopy, samples are analyzed by grinding with powdered potassium bromide and pressing into a disc. Alternatively, samples may be analyzed as a liquid film (neat, as a solution, or in a mull with Nujol) between two polished potassium bromide discs.
Due to its high solubility and hygroscopic nature it must be kept in a dry environment. The refractive index is about 1.55 at 1.0 μm.
Photography
In addition to manufacture of silver bromide, potassium bromide is used as a restrainer in black and white developer formulas. It improves differentiation between exposed and unexposed crystals of silver halide, and thus reduces fog.
| Physical sciences | Halide salts | Chemistry |
250515 | https://en.wikipedia.org/wiki/Landform | Landform | A landform is a natural or anthropogenic land feature on the solid surface of the Earth or other planetary body. Landforms together make up a given terrain, and their arrangement in the landscape is known as topography. Landforms include hills, mountains, canyons, and valleys, as well as shoreline features such as bays, peninsulas, and seas, including submerged features such as mid-ocean ridges, volcanoes, and the great ocean basins.
Physical characteristics
Landforms are categorized by characteristic physical attributes such as elevation, slope, orientation, structure
stratification, rock exposure, and soil type. Gross physical features or landforms include intuitive elements such as berms, mounds, hills, ridges, cliffs, valleys, rivers, peninsulas, volcanoes, and numerous other structural and size-scaled (e.g. ponds vs. lakes, hills vs. mountains) elements including various kinds of inland and oceanic waterbodies and sub-surface features. Mountains, hills, plateaux, and plains are the four major types of landforms. Minor landforms include buttes, canyons, valleys, and basins. Tectonic plate movement under the Earth can create landforms by pushing up mountains and hills.
Hierarchy of classes
Oceans and continents exemplify the highest-order landforms. Landform elements are parts of a high-order landforms that can be further identified and systematically given a cohesive definition such as hill-tops, shoulders, saddles, foreslopes and backslopes.
Some generic landform elements including: pits, peaks, channels, ridges, passes, pools and plains.
Terrain (or relief) is the third or vertical dimension of land surface. Topography is the study of terrain, although the word is often used as a synonym for relief itself. When relief is described underwater, the term bathymetry is used. In cartography, many different techniques are used to describe relief, including contour lines and triangulated irregular networks.
Elementary landforms (segments, facets, relief units) are the smallest homogeneous divisions of the land surface, at the given scale/resolution. These are areas with relatively homogeneous morphometric properties, bounded by lines of discontinuity. A plateau or a hill can be observed at various scales, ranging from a few hundred meters to hundreds of kilometers. Hence, the spatial distribution of landforms is often scale-dependent, as is the case for soils and geological strata.
A number of factors, ranging from plate tectonics to erosion and deposition (also due to human activity), can generate and affect landforms. Biological factors can also influence landforms—for example, note the role of vegetation in the development of dune systems and salt marshes, and the work of corals and algae in the formation of coral reefs.
Landforms do not include several man-made features, such as canals, ports and many harbors; and geographic features, such as deserts, forests, and grasslands. Many of the terms are not restricted to refer to features of the planet Earth, and can be used to describe surface features of other planets and similar objects in the Universe. Examples are mountains, hills, polar caps, and valleys, which are found on all of the terrestrial planets.
The scientific study of landforms is known as geomorphology.
In onomastic terminology, toponyms (geographical proper names) of individual landform objects (mountains, hills, valleys, etc.) are called oronyms.
Recent developments
Landforms may be extracted from a digital elevation model (DEM) using some automated techniques where the data has been gathered by modern satellites and stereoscopic aerial surveillance cameras. Until recently, compiling the data found in such data sets required time consuming and expensive techniques involving many man-hours. The most detailed DEMs available are measured directly using LIDAR techniques.
| Physical sciences | Landforms | null |
250540 | https://en.wikipedia.org/wiki/Rhodopsin | Rhodopsin | Rhodopsin, also known as visual purple, is a protein encoded by the RHO gene and a G-protein-coupled receptor (GPCR). It is a light-sensitive receptor protein that triggers visual phototransduction in rods. Rhodopsin mediates dim light vision and thus is extremely sensitive to light. When rhodopsin is exposed to light, it immediately photobleaches. In humans, it is regenerated fully in about 30 minutes, after which the rods are more sensitive. Defects in the rhodopsin gene cause eye diseases such as retinitis pigmentosa and congenital stationary night blindness.
Names
Rhodopsin was discovered by Franz Christian Boll in 1876. The name rhodopsin derives from Ancient Greek () for "rose", due to its pinkish color, and () for "sight". It was coined in 1878 by the German physiologist Wilhelm Friedrich Kühne (1837–1900).
When George Wald discovered that rhodopsin is a holoprotein, consisting of retinal and an apoprotein, he called it opsin, which today would be described more narrowly as apo-rhodopsin. Today, the term opsin refers more broadly to the class of G-protein-coupled receptors that bind retinal and as a result become a light sensitive photoreceptor, including all closely related proteins. When Wald and colleagues later isolated iodopsin from chicken retinas, thereby discovering the first known cone opsin, they called apo-iodopsin photopsin (for its relation to photopic vision) and apo-rhodopsin scotopsin (for its use in scotopic vision).
General
Rhodopsin is a protein found in the outer segment discs of rod cells. It mediates scotopic vision, which is monochromatic vision in dim light. Rhodopsin most strongly absorbs green-blue light (~500 nm) and appears therefore reddish-purple, hence the archaic term "visual purple".
Several closely related opsins differ only in a few amino acids and in the wavelengths of light that they absorb most strongly. Humans have, including rhodopsin, nine opsins, as well as cryptochrome (light-sensitive, but not an opsin).
Structure
Rhodopsin, like other opsins, is a G-protein-coupled receptor (GPCR). GPCRs are chemoreceptors that embed in the lipid bilayer of the cell membranes and have seven transmembrane domains forming a binding pocket for a ligand. The ligand for rhodopsin is the vitamin A-based chromophore 11-cis-retinal, which lies horizontally to the cell membrane and is covalently bound to a lysine residue (lys296) in the seventh transmembrane domain through a Schiff-base. However, 11-cis-retinal only blocks the binding pocket and does not activate rhodopsin. It is only activated when 11-cis-retinal absorbs a photon of light and isomerizes to all-trans-retinal, the receptor activating form, causing conformal changes in rhodopsin (bleaching), which activate a phototransduction cascade. Thus, a chemoreceptor is converted to a light or photo(n)receptor.
The retinal binding lysine is conserved in almost all opsins, only a few opsins having lost it during evolution. Opsins without the lysine are not light sensitive, including rhodopsin. Rhodopsin is made constitutively (continuously) active by some of those mutations even without light. Also wild-type rhodopsin is constitutively active, if no 11-cis-retinal is bound, but much less. Therefore 11-cis-retinal is an inverse agonist. Such mutations are one cause of autosomal dominant retinitis pigmentosa. Artificially, the retinal binding lysine can be shifted to other positions, even into other transmembrane domains, without changing the activity.
The rhodopsin of cattle has 348 amino acids, the retinal binding lysine being Lys296. It was the first opsin whose amino acid sequence and 3D-structure were determined. Its structure has been studied in detail by x-ray crystallography on rhodopsin crystals. Several models (e.g., the bicycle-pedal mechanism, hula-twist mechanism) attempt to explain how the retinal group can change its conformation without clashing with the enveloping rhodopsin protein pocket. Recent data support that rhodopsin is a functional monomer, instead of a dimer, which was the paradigm of G-protein-coupled receptors for many years.
Within its native membrane, rhodopsin is found at a high density facilitating its ability to capture photons. Due to its dense packing within the membrane, there is a higher chance of rhodopsin capturing proteins. However, the high density also provides a disadvantage when it comes to G protein signaling because the diffusion becomes more difficult in a crowded membrane that is packed with the receptor, rhodopsin.
Phototransduction
Rhodopsin is an essential G-protein coupled receptor in phototransduction.
Activation
In rhodopsin, the aldehyde group of retinal is covalently linked to the amino group of a lysine residue on the protein in a protonated Schiff base (-NH+=CH-). When rhodopsin absorbs light, its retinal cofactor isomerizes from the 11-cis to the all-trans configuration, and the protein subsequently undergoes a series of relaxations to accommodate the altered shape of the isomerized cofactor. The intermediates formed during this process were first investigated in the laboratory of George Wald, who received the Nobel prize for this research in 1967. The photoisomerization dynamics has been subsequently investigated with time-resolved IR spectroscopy and UV/Vis spectroscopy. A first photoproduct called photorhodopsin forms within 200 femtoseconds after irradiation, followed within picoseconds by a second one called bathorhodopsin with distorted all-trans bonds. This intermediate can be trapped and studied at cryogenic temperatures, and was initially referred to as prelumirhodopsin. In subsequent intermediates lumirhodopsin and metarhodopsin I, the Schiff's base linkage to all-trans retinal remains protonated, and the protein retains its reddish color. The critical change that initiates the neuronal excitation involves the conversion of metarhodopsin I to metarhodopsin II, which is associated with deprotonation of the Schiff's base and change in color from red to yellow.
Phototransduction cascade
The product of light activation, Metarhodopsin II, initiates the visual phototransduction second messenger pathway by stimulating the G-protein transducin (Gt), resulting in the liberation of its α subunit. This GTP-bound subunit in turn activates a cGMP phosphodiesterase. The cGMP phosphodiesterase hydrolyzes (breaks down) cGMP, lowering its local concentration so it can no longer activate cGMP-dependent cation channels. This leads to the hyperpolarization of photoreceptor cells, changing the rate at which they release transmitters.
Deactivation
Meta II (metarhodopsin II) is deactivated rapidly after activating transducin by rhodopsin kinase and arrestin. Rhodopsin pigment must be regenerated for further phototransduction to occur. This means replacing all-trans-retinal with 11-cis-retinal and the decay of Meta II is crucial in this process. During the decay of Meta II, the Schiff base link that normally holds all-trans-retinal and the apoprotein opsin (aporhodopsin) is hydrolyzed and becomes Meta III. In the rod outer segment, Meta III decays into separate all-trans-retinal and opsin. A second product of Meta II decay is an all-trans-retinal opsin complex in which the all-trans-retinal has been translocated to second binding sites. Whether the Meta II decay runs into Meta III or the all-trans-retinal opsin complex seems to depend on the pH of the reaction. Higher pH tends to drive the decay reaction towards Meta III.
Diseases of the retina
Mutations in the rhodopsin gene contribute majorly to various diseases of the retina such as retinitis pigmentosa. In general, the defect rhodopsin aggregates with ubiquitin in inclusion bodies, disrupts the intermediate filament network, and impairs the ability of the cell to degrade non-functioning proteins, which leads to photoreceptor apoptosis. Other mutations on rhodopsin lead to X-linked congenital stationary night blindness, mainly due to constitutive activation, when the mutations occur around the chromophore binding pocket of rhodopsin. Several other pathological states relating to rhodopsin have been discovered including poor post-Golgi trafficking, dysregulative activation, rod outer segment instability and arrestin binding.
| Biology and health sciences | Biological pigments | Biology |
250882 | https://en.wikipedia.org/wiki/Cat%27s%20Eye%20Nebula | Cat's Eye Nebula | The Cat's Eye Nebula (also known as NGC 6543 and Caldwell 6) is a planetary nebula in the northern constellation of Draco, discovered by William Herschel on February 15, 1786. It was the first planetary nebula whose spectrum was investigated by the English amateur astronomer William Huggins, demonstrating that planetary nebulae were gaseous and not stellar in nature. Structurally, the object has had high-resolution images by the Hubble Space Telescope revealing knots, jets, bubbles and complex arcs, being illuminated by the central hot planetary nebula nucleus (PNN).
It is a well-studied object that has been observed from radio to X-ray wavelengths. At the centre of the Cat's Eye Nebula is a dying Wolf Rayet star, the sort of which can be seen in the Webb Telescope's image of WR 124. The Cat's Eye Nebula's central star shines at magnitude +11.4. Hubble Space Telescope images show a sort of dart board pattern of concentric rings emanating outwards from the centre.
General information
NGC 6543 is a high northern declination deep-sky object. It has the combined magnitude of 8.1, with high surface brightness. Its small bright inner nebula subtends an average of 16.1 arcsec, with the outer prominent condensations about 25 arcsec. Deep images reveal an extended halo about 300 arcsec or 5 arcminutes across, that was once ejected by the central progenitor star during its red giant phase.
NGC 6543 is 4.4 minutes of arc from the current position of the north ecliptic pole, less than of the 45 arcminutes between Polaris and the current location of the Earth's northern axis of rotation. It is a convenient and accurate marker for the axis of rotation of the Earth's ecliptic, around which the celestial North Pole rotates. It is also a good marker for the nearby "invariable" axis of the solar system, which is the center of the circles which every planet's north pole, and the north pole of every planet's orbit, make in the sky. Since motion in the sky of the ecliptic pole is very slow compared to the motion of the Earth's north pole, its position as an ecliptic pole station marker is essentially permanent on the time-scale of human history, as opposed to the pole star, which changes every few thousand years.
Observations show the bright nebulosity has temperatures between and , whose densities average of about particles per cubic centimetre. Its outer halo has the higher temperature around , but is of much lower density. Velocity of the fast stellar wind is about , where spectroscopic analysis shows the current rate of mass loss averages solar masses per year, equivalent to twenty trillion tons per second (20 Eg/s).
Surface temperature for the central PNN is about , being 10,000 times as luminous as the sun. Stellar classification is O7 + [WR]-type star. Calculations suggest the PNN is over one solar mass, from a theoretical initial 5 solar masses. The central Wolf–Rayet star has a radius of (452,000 km). The Cat's Eye Nebula, given in some sources, lies about three thousand light-years from Earth.
Observations
The Cat's Eye was the first planetary nebula to be observed with a spectroscope by William Huggins on August 29, 1864. Huggins' observations revealed that the nebula's spectrum was non-continuous and made of a few bright emission lines, first indication that planetary nebulae consist of tenuous ionised gas. Spectroscopic observations at these wavelengths are used in abundance determinations, while images at these wavelengths have been used to reveal the intricate structure of the nebula.
Infrared observations
Observations of NGC 6543 at far-infrared wavelengths (about 60 μm) reveal the presence of stellar dust at low temperatures. The dust is believed to have formed during the last phases of the progenitor star's life. It absorbs light from the central star and re-radiates it at infrared wavelengths. The spectrum of the infrared dust emission implies that the dust temperature is about 85 K, while the mass of the dust is estimated at solar masses.
Infrared emission also reveals the presence of un-ionised material such as molecular hydrogen (H2) and argon. In many planetary nebulae, molecular emission is greatest at larger distances from the star, where more material is un-ionised, but molecular hydrogen emission in NGC 6543 seems to be bright at the inner edge of its outer halo. This may be due to shock waves exciting the H2 as ejecta moving at different speeds collide. The overall appearance of the Cat's Eye Nebula in infrared (wavelengths 2–8 μm) is similar in visible light.
Optical and ultraviolet observations
The Hubble Space Telescope image produced here is in false colour, designed to highlight regions of high and low ionisation. Three images were taken, in filters isolating the light emitted by singly ionised hydrogen at 656.3 nm, singly ionised nitrogen at 658.4 nm and doubly ionised oxygen at 500.7 nm. The images were combined as red, green and blue channels respectively, although their true colours are red, red and green. The image reveals two "caps" of less ionised material at the edge of the nebula.
X-ray observations
In 2001, observations at X-ray wavelengths by the Chandra X-ray Observatory revealed the presence of extremely hot gas within NGC 6543 with the temperature of . It is thought that the very hot gas results from the violent interaction of a fast stellar wind with material previously ejected. This interaction has hollowed out the inner bubble of the nebula. Chandra observations have also revealed a point source at the position of the central star. The spectrum of this source extends to the hard part of the X-ray spectrum, to 0.5–. A star with the photospheric temperature of about would not be expected to emit strongly in hard X-rays, and so their presence is something of a mystery. It may suggest the presence of a high temperature accretion disk within a binary star system. The hard X-ray data remain intriguing more than ten years later: the Cat's Eye was included in a 2012 Chandra survey of 21 central stars of planetary nebulae (CSPNe) in the solar neighborhood, which found: "All but one of the X-ray point sources detected at CSPNe display X-ray spectra that are harder than expected from hot (~) central star photospheres, possibly indicating a high frequency of binary companions to CSPNe. Other potential explanations include self-shocking winds or PN mass fallback."
Distance
Planetary nebulae distances like NGC 6543 are generally very inaccurate and not well known. Some recent Hubble Space Telescope observations of NGC 6543 taken several years apart determine its distance from the angular expansion rate of 3.457 milliarcseconds per year. Assuming a line of sight expansion velocity of 16.4 km·s−1, this implies that NGC 6543's distance is parsecs ( or light-years) away from Earth. Several other distance references, like what is quoted in SIMBAD in 2014 based on Stanghellini, L., et al. (2008) suggest the distance is parsecs ( light-years).
Age
The angular expansion of the nebula can also be used to estimate its age. If it has been expanding at a constant rate of 10 milliarcseconds a year, then it would take to reach a diameter of 20 arcseconds. This may be an upper limit to the age, because ejected material will be slowed when it encounters material ejected from the star at earlier stages of its evolution, and the interstellar medium.
Composition
Like most astronomical objects, NGC 6543 consists mostly of hydrogen and helium, with heavier elements present in small quantities. The exact composition may be determined by spectroscopic studies. Abundances are generally expressed relative to hydrogen, the most abundant element.
Different studies generally find varying values for elemental abundances. This is often because spectrographs attached to telescopes do not collect all the light from objects being observed, instead gathering light from a slit or small aperture. Therefore, different observations may sample different parts of the nebula.
However, results for NGC 6543 broadly agree that, relative to hydrogen, the helium abundance is about 0.12, carbon and nitrogen abundances are both about , and the oxygen abundance is about . These are fairly typical abundances for planetary nebulae, with the carbon, nitrogen and oxygen abundances all larger than the values found for the sun, due to the effects of nucleosynthesis enriching the star's atmosphere in heavy elements before it is ejected as a planetary nebula.
Deep spectroscopic analysis of NGC 6543 may indicate that the nebula contains a small amount of material which is highly enriched in heavy elements; this is discussed below.
Kinematics and morphology
The Cat's Eye Nebula is structurally a very complex nebula, and the mechanism or mechanisms that have given rise to its complicated morphology are not well understood. The central bright part of the nebula consists of the inner elongated bubble (inner ellipse) filled with hot gas. It, in turn, is nested into a pair of larger spherical bubbles conjoined together along their waist. The waist is observed as the second larger ellipse lying perpendicular to the bubble with hot gas.
The structure of the bright portion of the nebula is primarily caused by the interaction of a fast stellar wind being emitted by the central PNN with the visible material ejected during the formation of the nebula. This interaction causes the emission of X-rays discussed above. The stellar wind, blowing with the velocity as high as , has 'hollowed out' the inner bubble of the nebula, and appears to have burst the bubble at both ends.
It is also suspected that the central WR:+O7 spectral class PNN star, HD 164963 / BD +66 1066 / PPM 20679 of the nebula may be generated by a binary star. The existence of an accretion disk caused by mass transfer between the two components of the system may give rise to astronomical jets, which would interact with previously ejected material. Over time, the direction of the jets would vary due to precession.
Outside the bright inner portion of the nebula, there are a series of concentric rings, thought to have been ejected before the formation of the planetary nebula, while the star was on the asymptotic giant branch of the Hertzsprung–Russell diagram. These rings are very evenly spaced, suggesting that the mechanism responsible for their formation ejected them at very regular intervals and at very similar speeds. The total mass of the rings is about 0.1 solar masses. The pulsations that formed the rings probably started 15,000 years ago and ceased about years ago, when the formation of the bright central part began (see above).
Further, a large faint halo extends to large distances from the star. The halo again predates the formation of the main nebula. The mass of the halo is estimated as 0.26–0.92 solar masses.
| Physical sciences | Notable nebulae | Astronomy |
250896 | https://en.wikipedia.org/wiki/Implicit%20function | Implicit function | In mathematics, an implicit equation is a relation of the form where is a function of several variables (often a polynomial). For example, the implicit equation of the unit circle is
An implicit function is a function that is defined by an implicit equation, that relates one of the variables, considered as the value of the function, with the others considered as the arguments. For example, the equation of the unit circle defines as an implicit function of if , and is restricted to nonnegative values.
The implicit function theorem provides conditions under which some kinds of implicit equations define implicit functions, namely those that are obtained by equating to zero multivariable functions that are continuously differentiable.
Examples
Inverse functions
A common type of implicit function is an inverse function. Not all functions have a unique inverse function. If is a function of that has a unique inverse, then the inverse function of , called , is the unique function giving a solution of the equation
for in terms of . This solution can then be written as
Defining as the inverse of is an implicit definition. For some functions , can be written out explicitly as a closed-form expression — for instance, if , then . However, this is often not possible, or only by introducing a new notation (as in the product log example below).
Intuitively, an inverse function is obtained from by interchanging the roles of the dependent and independent variables.
Example: The product log is an implicit function giving the solution for of the equation .
Algebraic functions
An algebraic function is a function that satisfies a polynomial equation whose coefficients are themselves polynomials. For example, an algebraic function in one variable gives a solution for of an equation
where the coefficients are polynomial functions of . This algebraic function can be written as the right side of the solution equation . Written like this, is a multi-valued implicit function.
Algebraic functions play an important role in mathematical analysis and algebraic geometry. A simple example of an algebraic function is given by the left side of the unit circle equation:
Solving for gives an explicit solution:
But even without specifying this explicit solution, it is possible to refer to the implicit solution of the unit circle equation as , where is the multi-valued implicit function.
While explicit solutions can be found for equations that are quadratic, cubic, and quartic in , the same is not in general true for quintic and higher degree equations, such as
Nevertheless, one can still refer to the implicit solution involving the multi-valued implicit function .
Caveats
Not every equation implies a graph of a single-valued function, the circle equation being one prominent example. Another example is an implicit function given by where is a cubic polynomial having a "hump" in its graph. Thus, for an implicit function to be a true (single-valued) function it might be necessary to use just part of the graph. An implicit function can sometimes be successfully defined as a true function only after "zooming in" on some part of the -axis and "cutting away" some unwanted function branches. Then an equation expressing as an implicit function of the other variables can be written.
The defining equation can also have other pathologies. For example, the equation does not imply a function giving solutions for at all; it is a vertical line. In order to avoid a problem like this, various constraints are frequently imposed on the allowable sorts of equations or on the domain. The implicit function theorem provides a uniform way of handling these sorts of pathologies.
Implicit differentiation
In calculus, a method called implicit differentiation makes use of the chain rule to differentiate implicitly defined functions.
To differentiate an implicit function , defined by an equation , it is not generally possible to solve it explicitly for and then differentiate. Instead, one can totally differentiate with respect to and and then solve the resulting linear equation for to explicitly get the derivative in terms of and . Even when it is possible to explicitly solve the original equation, the formula resulting from total differentiation is, in general, much simpler and easier to use.
Examples
Example 1
Consider
This equation is easy to solve for , giving
where the right side is the explicit form of the function . Differentiation then gives .
Alternatively, one can totally differentiate the original equation:
Solving for gives
the same answer as obtained previously.
Example 2
An example of an implicit function for which implicit differentiation is easier than using explicit differentiation is the function defined by the equation
To differentiate this explicitly with respect to , one has first to get
and then differentiate this function. This creates two derivatives: one for and another for .
It is substantially easier to implicitly differentiate the original equation:
giving
Example 3
Often, it is difficult or impossible to solve explicitly for , and implicit differentiation is the only feasible method of differentiation. An example is the equation
It is impossible to algebraically express explicitly as a function of , and therefore one cannot find by explicit differentiation. Using the implicit method, can be obtained by differentiating the equation to obtain
where . Factoring out shows that
which yields the result
which is defined for
General formula for derivative of implicit function
If , the derivative of the implicit function is given by
where and indicate the partial derivatives of with respect to and .
The above formula comes from using the generalized chain rule to obtain the total derivative — with respect to — of both sides of :
hence
which, when solved for , gives the expression above.
Implicit function theorem
Let be a differentiable function of two variables, and be a pair of real numbers such that . If , then defines an implicit function that is differentiable in some small enough neighbourhood of ; in other words, there is a differentiable function that is defined and differentiable in some neighbourhood of , such that for in this neighbourhood.
The condition means that is a regular point of the implicit curve of implicit equation where the tangent is not vertical.
In a less technical language, implicit functions exist and can be differentiated, if the curve has a non-vertical tangent.
In algebraic geometry
Consider a relation of the form , where is a multivariable polynomial. The set of the values of the variables that satisfy this relation is called an implicit curve if and an implicit surface if . The implicit equations are the basis of algebraic geometry, whose basic subjects of study are the simultaneous solutions of several implicit equations whose left-hand sides are polynomials. These sets of simultaneous solutions are called affine algebraic sets.
In differential equations
The solutions of differential equations generally appear expressed by an implicit function.
Applications in economics
Marginal rate of substitution
In economics, when the level set is an indifference curve for the quantities and consumed of two goods, the absolute value of the implicit derivative is interpreted as the marginal rate of substitution of the two goods: how much more of one must receive in order to be indifferent to a loss of one unit of .
Marginal rate of technical substitution
Similarly, sometimes the level set is an isoquant showing various combinations of utilized quantities of labor and of physical capital each of which would result in the production of the same given quantity of output of some good. In this case the absolute value of the implicit derivative is interpreted as the marginal rate of technical substitution between the two factors of production: how much more capital the firm must use to produce the same amount of output with one less unit of labor.
Optimization
Often in economic theory, some function such as a utility function or a profit function is to be maximized with respect to a choice vector even though the objective function has not been restricted to any specific functional form. The implicit function theorem guarantees that the first-order conditions of the optimization define an implicit function for each element of the optimal vector of the choice vector . When profit is being maximized, typically the resulting implicit functions are the labor demand function and the supply functions of various goods. When utility is being maximized, typically the resulting implicit functions are the labor supply function and the demand functions for various goods.
Moreover, the influence of the problem's parameters on — the partial derivatives of the implicit function — can be expressed as total derivatives of the system of first-order conditions found using total differentiation.
| Mathematics | Basics | null |
250959 | https://en.wikipedia.org/wiki/Laramide%20orogeny | Laramide orogeny | The Laramide orogeny was a time period of mountain building in western North America, which started in the Late Cretaceous, 80 to 70 million years ago, and ended 55 to 35 million years ago. The exact duration and ages of beginning and end of the orogeny are in dispute. The Laramide orogeny occurred in a series of pulses, with quiescent phases intervening. The major feature that was created by this orogeny was deep-seated, thick-skinned deformation, with evidence of this orogeny found from Canada to northern Mexico, with the easternmost extent of the mountain-building represented by the Black Hills of South Dakota. The phenomenon is named for the Laramie Mountains of eastern Wyoming. The Laramide orogeny is sometimes confused with the Sevier orogeny, which partially overlapped in time and space.
The orogeny is commonly attributed to events off the west coast of North America, where the Kula and Farallon Plates were sliding under the North American Plate. Most hypotheses propose that oceanic crust was undergoing flat-slab subduction, that is, subduction at a shallow angle. As a consequence, no magmatism occurred in the central west of the continent, and the underlying oceanic lithosphere actually caused drag on the root of the overlying continental lithosphere. One cause for shallow subduction may have been an increased rate of plate convergence. Another proposed cause was subduction of thickened oceanic crust.
Magmatism associated with subduction occurred not near the plate edges (as in the volcanic arc of the Andes, for example), but far to the east, along the Colorado Mineral Belt. Geologists call such a lack of volcanic activity near a subduction zone a magmatic gap. This particular gap may have occurred because the subducted slab was in contact with relatively cool continental lithosphere, not hotter asthenosphere. One result of shallow angle of subduction and the drag that it caused was a broad belt of mountains, some of which were the progenitors of the Rocky Mountains. Part of the proto-Rocky Mountains would be later modified by extension to become the Basin and Range Province.
Basins and mountains
The Laramide orogeny produced intermontane structural basins and adjacent mountain blocks by means of deformation. This style of deformation is typical of continental plates adjacent to convergent margins of long duration that have not sustained continent/continent collisions. This tectonic setting produces a pattern of compressive uplifts and basins, with most of the deformation confined to block edges. Twelve kilometers of structural relief between basins and adjacent uplifts is not uncommon. The basins contain several thousand meters of Paleozoic and Mesozoic sedimentary rocks that predate the Laramide orogeny. As much as of Cretaceous and Cenozoic sediments filled these orogenically-defined basins. Deformed Paleocene and Eocene deposits record continuing orogenic activity.
During the Laramide orogeny, basin floors and mountain summits were much closer to sea level than today. After the seas retreated from the Rocky Mountain region, floodplains, swamps, and vast lakes developed in the basins. Drainage systems imposed at that time persist today. Since the Oligocene, episodic epeirogenic uplift gradually raised the entire region, including the Great Plains, to present elevations. Most of the modern topography is the result of Pliocene and Pleistocene events, including additional uplift, glaciation of the high country, and denudation and dissection of older Cenozoic surfaces in the basin by fluvial processes.
In the United States, these distinctive intermontane basins occur principally in the central Rocky Mountains from Colorado and Utah (Uinta Basin) to Montana and are best developed in Wyoming, with the Bighorn, Powder River, and Wind River being the largest. Topographically, the basin floors resemble the surface of the western Great Plains, except for vistas of surrounding mountains.
At most boundaries, Paleozoic through Paleogene units dip steeply into the basins off uplifted blocks cored by Precambrian rocks. The eroded steeply dipping units form hogbacks and flatirons. Many of the boundaries are thrust or reverse faults. Although other boundaries appear to be monoclinal flexures, faulting is suspected at depth. Most bounding faults show evidence of at least two episodes of Laramide (Late Cretaceous and Eocene) movement, suggesting both thrust and strike-slip types of displacement.
Ecological consequences
According to paleontologist Thomas M. Lehman, the Laramide orogeny triggered "the most dramatic event that affected Late Cretaceous dinosaur communities in North America prior to their extinction." This turnover event saw the replacement of specialized and highly ornamented centrosaurine and lambeosaurines by more basal upland dinosaurs in the south, while northern biomes became dominated by Triceratops with a greatly reduced hadrosaur community.
| Physical sciences | Geologic features | Earth science |
250984 | https://en.wikipedia.org/wiki/Divergent%20boundary | Divergent boundary | In plate tectonics, a divergent boundary or divergent plate boundary (also known as a constructive boundary or an extensional boundary) is a linear feature that exists between two tectonic plates that are moving away from each other. Divergent boundaries within continents initially produce rifts, which eventually become rift valleys. Most active divergent plate boundaries occur between oceanic plates and exist as mid-oceanic ridges.
Current research indicates that complex convection within the Earth's mantle allows material to rise to the base of the lithosphere beneath each divergent plate boundary.
This supplies the area with huge amounts of heat and a reduction in pressure that melts rock from the asthenosphere (or upper mantle) beneath the rift area, forming large flood basalt or lava flows. Each eruption occurs in only a part of the plate boundary at any one time, but when it does occur, it fills in the opening gap as the two opposing plates move away from each other.
Over millions of years, tectonic plates may move many hundreds of kilometers away from both sides of a divergent plate boundary. Because of this, rocks closest to a boundary are younger than rocks further away on the same plate.
Description
At divergent boundaries, two plates move away from each other and the space that this creates is filled with new crustal material sourced from molten magma that forms below. The origin of new divergent boundaries at triple junctions is sometimes thought to be associated with the phenomenon known as hotspots. Here, exceedingly large convective cells bring very large quantities of hot asthenospheric material near the surface, and the kinetic energy is thought to be sufficient to break apart the lithosphere.
Divergent boundaries are typified in the oceanic lithosphere by the rifts of the oceanic ridge system, including the Mid-Atlantic Ridge and the East Pacific Rise, and in the continental lithosphere by rift valleys such as the famous East African Great Rift Valley. Divergent boundaries can create massive fault zones in the oceanic ridge system. Spreading is generally not uniform, so where spreading rates of adjacent ridge blocks are different, massive transform faults occur. These are the fracture zones, many bearing names, that are a major source of submarine earthquakes. A seafloor map will show a rather strange pattern of blocky structures that are separated by linear features perpendicular to the ridge axis. If one views the seafloor between the fracture zones as conveyor belts carrying the ridge on each side of the rift away from the spreading center the action becomes clear. Crest depths of the old ridges, parallel to the current spreading center, will be older and deeper... (from thermal contraction and subsidence).
It is at mid-ocean ridges that one of the key pieces of evidence forcing acceptance of the seafloor spreading hypothesis was found. Airborne geomagnetic surveys showed a strange pattern of symmetrical magnetic reversals on opposite sides of ridge centers. The pattern was far too regular to be coincidental as the widths of the opposing bands were too closely matched. Scientists had been studying polar reversals and the link was made by Lawrence W. Morley, Frederick John Vine and Drummond Hoyle Matthews in the Morley–Vine–Matthews hypothesis. The magnetic banding directly corresponds with the Earth's polar reversals. This was confirmed by measuring the ages of the rocks within each band. The banding furnishes a map in time and space of both spreading rate and polar reversals.
Examples
Mid-Atlantic Ridge
Red Sea Rift
Baikal Rift Zone – incipient plate boundary
East African Rift – incipient plate boundary
East Pacific Rise
Gakkel Ridge
Cocos–Nazca spreading centre
Explorer Ridge
Juan de Fuca Ridge
Pacific-Antarctic Ridge
Southeast Indian Ridge
Other plate boundary types
Convergent boundary
Transform boundary
| Physical sciences | Tectonics | Earth science |
251003 | https://en.wikipedia.org/wiki/Hyoid%20bone | Hyoid bone | The hyoid bone (lingual bone or tongue-bone) () is a horseshoe-shaped bone situated in the anterior midline of the neck between the chin and the thyroid cartilage. At rest, it lies between the base of the mandible and the third cervical vertebra.
Unlike other bones, the hyoid is only distantly articulated to other bones by muscles or ligaments. It is the only bone in the human body that is not connected to any other bones. The hyoid is anchored by muscles from the anterior, posterior and inferior directions, and aids in tongue movement and swallowing. The hyoid bone provides attachment to the muscles of the floor of the mouth and the tongue above, the larynx below, and the epiglottis and pharynx behind.
Its name is derived .
Structure
The hyoid bone is classed as an irregular bone and consists of a central part called the body, and two pairs of horns, the greater and lesser horns.
Body
The body of the hyoid bone is the central part of the hyoid bone.
At the front, the body is convex and directed forward and upward.
It is crossed in its upper half by a well-marked transverse ridge with a slight downward convexity, and in many cases a vertical median ridge divides it into two lateral halves.
The portion of the vertical ridge above the transverse line is present in a majority of specimens, but the lower portion is evident only in rare cases.
The anterior surface gives insertion to the geniohyoid muscle in the greater part of its extent both above and below the transverse ridge; a portion of the origin of the hyoglossus notches the lateral margin of the geniohyoid attachment.
Below the transverse ridge the mylohyoid, sternohyoid, and omohyoid are inserted.
At the back, the smooth, concave, directed backward and downward, and separated from the epiglottis by the hyothyroid membrane and a quantity of loose areolar tissue; a bursa intervenes between it and the hyothyroid membrane.
Above, the body is rounded, and gives attachment to the hyothyroid membrane and some aponeurotic fibers of the genioglossus.
Below, the body affords insertion medially to the sternohyoid and laterally to the omohyoid and occasionally a portion of the thyrohyoid. It also gives attachment to the Levator glandulae thyreoideae, when this muscle is present.
Horns
The greater and lesser horns () are two sections of bone that project from each side of the hyoid.
Greater horns
The greater horns project backward from the outer borders of the body; they are flattened from above downward and taper to their end, which is a bony tubercle connecting to the lateral thyrohyoid ligament. The upper surface of the greater horns are rough and close to its lateral border, and facilitates muscular attachment. The largest of muscles that attach to the upper surface of the greater horns are the hyoglossus and the middle pharyngeal constrictor, which extend along the whole length of the horns; the digastric muscle and stylohyoid muscle have small insertions in front of these near the junction of the body with the horns. To the medial border, the thyrohyoid membrane is attached, while the anterior half of the lateral border gives insertion to the thyrohyoid muscle.
Lesser horns
The lesser horns are two small, conical eminences, attached by their bases to the angles of junction between the body and greater horns of the hyoid bone. They are connected to the body of the bone by fibrous tissue, and occasionally to the greater horns by distinct diarthrodial joints, which usually persist throughout life, but occasionally become ankylosed. The lesser horns are situated in the line of the transverse ridge on the body and appear to be continuations of it. The apex of each horn gives attachment to the stylohyoid ligament; the chondroglossus rises from the medial side of the base.
Development
The second pharyngeal arch, also called the hyoid arch, gives rise to the lesser cornu of the hyoid and the upper part of the body of the hyoid. The cartilage of the third pharyngeal arch forms the greater cornu of the hyoid and the lower portion of the body of the hyoid.
The hyoid is ossified from six centers: two for the body, and one for each cornu. Ossification commences in the greater cornua toward the end of fetal development, in the hyoid body shortly afterward, and in the lesser cornua during the first or second year after birth. Until middle age, the connection between the body and greater cornu is fibrous.
In early life, the outer borders of the body are connected to the greater horns by synchondroses; after middle life, usually by bony union.
Blood supply
Blood is supplied to the hyoid bone via the lingual artery, which runs down from the tongue to the greater horns of the bone. The suprahyoid branch of the lingual artery runs along the upper border of the hyoid bone and supplies blood to the attached muscles.
Function
The hyoid bone is present in many mammals. It allows a wider range of tongue, pharyngeal and laryngeal movements by bracing these structures alongside each other in order to produce variation. Its descent in living creatures is not unique to Homo sapiens, and does not allow the production of a wide range of sounds: with a lower larynx, men do not produce a wider range of sounds than women and two-year-old babies. Moreover, the larynx position of Neanderthals was not a handicap to producing speech sounds. The discovery of a modern-looking hyoid bone of a Neanderthal man in the Kebara Cave in Israel led its discoverers to argue that the Neanderthals had a descended larynx, and thus human-like speech capabilities. However, other researchers have claimed that the morphology of the hyoid is not indicative of the larynx's position. Recent research has indicated that the hyoid bone may have significant involvement in the ability to swallow. It has been hypothesized that the mammalian hyoid bone evolved in conjunction with the development of lactation, thus allowing babies to suckle milk. It is necessary to take into consideration the skull base, the mandible and the cervical vertebrae and a cranial reference plane.
Muscle attachments
A large number of muscles attach to the hyoid:
Superior
Middle pharyngeal constrictor muscle
Hyoglossus muscle
Genioglossus
Intrinsic muscles of the tongue
Suprahyoid muscles
Digastric muscle
Stylohyoid muscle
Geniohyoid muscle
Mylohyoid muscle
Inferior
Thyrohyoid muscle
Omohyoid muscle
Sternohyoid muscle
Sternothyroid muscle
Clinical significance
The hyoid bone is important to a number of physiological functions, including breathing, swallowing and speech. It is also thought to play a key role in keeping the upper airway open during sleep, and as such, the development and treatment of obstructive sleep apnea (OSA; characterized by repetitive collapse of the upper airway during sleep). A mechanistic involvement of the hyoid bone in OSA is supported by numerous studies demonstrating that a more inferiorly positioned hyoid bone is strongly associated with the presence and severity of the disorder. Movement of the hyoid bone is also thought to be important in modifying upper airway properties, which was recently demonstrated in computer model simulations. A surgical procedure that aims to potentially increase and improve the airway is called hyoid suspension.
Due to its position, the hyoid bone is not easily susceptible to fracture. In a suspected case of murder or physical abuse of an adult, a fractured hyoid strongly indicates throttling or strangulation. In children and adolescents (in whom the hyoid bone is still flexible because ossification is yet to be completed) fracture may not occur even after serious trauma.
Other animals
The hyoid bone is derived from the lower half of the second gill arch in fish, which separates the first gill slit from the spiracle, and is often called the hyoid arch. In many vertebrates, it also incorporates elements of other gill arches, and has a correspondingly greater number of cornua. Amphibians and non-avian reptiles may have many cornua, while mammals (including humans) have two pairs, and birds only one. In birds, and some reptiles, the body of the hyoid is greatly extended forward, creating a solid bony support for the tongue. The howler monkey Alouatta has a pneumatized hyoid bone, one of the few cases of postcranial pneumatization of bones outside Saurischia.
In woodpeckers, the hyoid bone is elongated, with the horns wrapping around the back of the skull. This is part of the system that keeps the brain cushioned and undamaged by the pecking action.
In mammals, the hyoid often determines whether one can roar. If the hyoid is incompletely ossified (for example: lions), it allows the animal to roar, but not purr. If the hyoid is completely ossified (for example: cheetahs), it does not allow the animal to roar, but instead will allow the animal to purr and meow, as seen in house cats (lions, cheetahs and house cats all belong to the family Felidae).
In veterinary anatomy, the term hyoid apparatus is the collective term used to refer to the bones of the tongue—a pair of stylohyoidea, a pair of thyrohyoidea, and unpaired basihyoideum—and associated, upper-gular connective tissues. In humans, the single hyoid bone is an equivalent of the hyoid apparatus.
| Biology and health sciences | Skeletal system | Biology |
251010 | https://en.wikipedia.org/wiki/Whiskers | Whiskers | Whiskers or vibrissae (; : vibrissa; ) are a type of stiff, functional hair used by most therian mammals to sense their environment. These hairs are finely specialised for this purpose, whereas other types of hair are coarser as tactile sensors. Although whiskers are specifically those found around the face, vibrissae are known to grow in clusters at various places around the body. Most mammals have them, including all non-human primates and especially nocturnal mammals. Monotremes, however, lack them.
Whiskers are sensitive tactile hairs that aid navigation, locomotion, exploration, hunting, social touch and perform other functions.
This article is primarily about the specialised sensing hairs of mammals, but some birds, fish, insects, crustaceans and other arthropods are known to have similar structures also used to sense the environment.
Etymology
Vibrissae (from Latin 'to vibrate') from the characteristic motion seen in a small rodent that is otherwise sitting still.
In medicine, the term also refers to the thick hairs found inside human nostrils.
Evolution
The last common ancestor of all extant mammals had vibrissae. All other extant mammal species besides great apes retain the same ancestral layout of the whiskers along with the special facial muscles that move them.
Anatomy
Vibrissae are anatomically distinguished from other hair. They are easily visually identified since they are longer, stiffer, significantly larger in diameter, and stand above the surrounding fur by a considerable amount. In addition, they have well-innervated follicles, and an identifiable representation in the somatosensory cortex of the brain. The largest number and the longest are found among the small, social, arboreal, and nocturnal mammals. Whiskers of aquatic mammals are the most sensitive. During foraging in complex, dark habitats, whiskers are rapidly moved in a cyclic way, tracing small circles at their tips. This motion, called "whisking" can occur at speeds of 25 Hz in mice, which is one of the fastest movements that mammals can make. Small animals use whisking to position their front paws during locomotion.
Vibrissal groups
Vibrissae typically grow in clusters. These groups vary somewhat in form and function, but they are relatively consistent among land mammals. Between land and marine mammals, there is less consistency (though commonalities are certainly present).
Many land mammals, like rats and hamsters, have four typical whisker groups on their heads (called cranial vibrissae), which might vary among animals due to different lifestyles. These cranial groups include:
above the eyes (supraorbital)
on the cheeks (genal)
where a moustache would be (mystacial)
under the snout (mandibular).
The mystacial whiskers can be roughly identified as macrovibrissae (long whiskers for feeling the space around the head) and microvibrissae (small, down-pointing whiskers for identifying objects). Not only are these two types hard to distinguish on an animal's face (see for example the image of a rat here), there are similarly weak distinctions in how they are used, though the distinction is nonetheless referred to ubiquitously in scientific literature and is considered useful in analysis.
Many land mammals, including domestic cats, also have vibrissae on the underside of the leg just above the paws (called carpal vibrissae). Whilst these five major groups are often reported in studies of land mammals, several other groups have been reported more occasionally; for instance nasal, angular, and submental whiskers.
Marine mammals can have substantially different arrangements of their vibrissae. For instance, whales and dolphins have lost their snout whiskers and gained vibrissae around their blowholes, whereas every single one of the body hairs of the Florida manatee may be a vibrissa (see image). Other marine mammals, like seals and sea-lions, have head vibrissae just like those on land mammals (see image), although these groups function quite differently.
Vibrissal follicles have evolved other functions in dolphins, such as electroreception.
Vibrissae
The vibrissal hair is usually thicker and stiffer than other types of (pelagic) hair but, like other hairs, the shaft consists of an inert material (keratin) and contains no nerves. However, vibrissae are different from other hair structures because they grow from a special hair follicle incorporating a capsule of blood called a blood sinus which is heavily innervated by sensory nerves. Vibrissae are symmetrically arranged in groups on the face and supply the trigeminal nerve.
The mystacial macrovibrissae are shared by a large group of land and marine mammals (see images), and it is this group that has received by far the most scientific study. The arrangement of these whiskers is not random: they form an ordered grid of arcs (columns) and rows, with shorter whiskers at the front and longer whiskers at the rear (see images). In the mouse, gerbil, hamster, rat, guinea pig, rabbit, and cat, each individual follicle is innervated by 100–200 primary afferent nerve cells. These cells serve an even larger number of mechanoreceptors of at least eight distinct types. Accordingly, even small deflections of the vibrissal hair can evoke a sensory response in the animal. Rats and mice typically have approximately 30 macrovibrissae on each side of the face, with whisker lengths up to around 50 mm in (laboratory) rats, 30 mm in (laboratory) mice, and a slightly larger number of microvibrissae. Thus, an estimate for the total number of sensory nerve cells serving the mystacial vibrissal array on the face of a rat or mouse might be 25,000. Natural shapes of rat's mystacial pad vibrissae are well approximated by pieces of the Euler spiral. When all these pieces for a single rat are assembled together, they span an interval extending from one coiled domain of the Euler spiral to the other.
Marine mammals may make even greater investment in their vibrissal sensory system than rats and mice. Seal whiskers, which are similarly arrayed across the mystacial region, are each served by around 10 times as many nerve fibres as those in rats and mice, so that the total number of nerve cells innervating the mystacial vibrissae of a seal has been estimated to be in excess of 300,000. Manatees, remarkably, have around 600 vibrissae on or around their lips.
Whiskers can be very long in some species; the length of a chinchilla's whiskers can be more than a third of its body length (see image). Even in species with shorter whiskers, they can be very prominent appendages (see images). Thus, whilst whiskers certainly could be described as "proximal sensors" in contrast to, say, eyes, they offer a tactile sense with a sensing range that is functionally very significant.
Operation
Movement
The follicles of some groups of vibrissae in some species are motile. Generally, the supraorbital, genal and macrovibrissae are motile, whereas the microvibrissae are not. This is reflected in anatomical reports that have identified musculature associated with the macrovibrissae that is absent for the microvibrissae. A small muscle 'sling' is attached to each macrovibrissa and can move it more-or-less independently of the others, whilst larger muscles in the surrounding tissue move many or all of the macrovibrissae together.
Amongst those species with motile macrovibrissae, some (rats, mice, flying squirrels, gerbils, chinchillas, hamsters, shrews, porcupines, opossums) move them back and forth periodically in a movement known as whisking, while other species (cats, dogs, raccoons, pandas) do not appear to. The distribution of mechanoreceptor types in the whisker follicle differs between rats and cats, which may correspond to this difference in the way they are used. Whisking movements are amongst the fastest produced by mammals. In all whisking animals in which it has so far been measured, these whisking movements are rapidly controlled in response to behavioural and environmental conditions. The whisking movements occur in bouts of variable duration, and at rates between 3 and 25 whisks/second. Movements of the whiskers are closely coordinated with those of the head and body.
Function
Generally, vibrissae are considered to mediate a tactile sense, complementary to that of skin. This is presumed to be advantageous in particular to animals that cannot always rely on sight to navigate or to find food, for example, nocturnal animals or animals which forage in muddy waters. Whiskers can also function as wind detecting antannae such as the supra-orbital ones in rats.
Sensory function aside, movements of the vibrissae may also indicate something of the state of mind of the animal, and the whiskers play a role in social behaviour of rats.
The sensory function of vibrissae is an active research area—experiments to establish the capabilities of whiskers use a variety of techniques, including temporary deprivation either of the whisker sense or of other senses. Animals can be deprived of their whisker sense for a period of weeks by whisker trimming (they soon grow back), or for the duration of an experimental trial by restraining the whiskers with a flexible cover like a mask (the latter technique is used, in particular, in studies of marine mammals). Such experiments have shown that whiskers are required for, or contribute to: object localization, orienting of the snout, detection of movement, texture discrimination, shape discrimination, exploration, thigmotaxis, locomotion, maintenance of equilibrium, maze learning, swimming, locating food pellets, locating food animals, and fighting, as well as nipple attachment and huddling in rat pups.
Whisking—the periodic movement of the whiskers—is also presumed to serve tactile sensing in some way. However, exactly why an animal might be driven "to beat the night with sticks", as one researcher once put it, is a matter of debate, and the answer is probably multi-faceted. Scholarpedia offers:
Animals that do not whisk, but have motile whiskers, presumably also gain some advantage from the investment in musculature. Dorothy Souza, in her book Look What Whiskers Can Do reports some whisker movement during prey capture (in cats, in this case):
Anecdotally, it is often stated that cats use their whiskers to gauge whether an opening is wide enough for their body to pass through. This is sometimes supported by the statement that the whiskers of individual cats extend out to about the same width as the cat's body, but at least two informal reports indicate that whisker length is genetically determined and does not vary as the cat grows thinner or fatter. In the laboratory, rats are able to accurately (within 5–10%) discriminate the size of an opening, so it seems likely that cats can use their whiskers for this purpose. However, reports of cats, particularly kittens, with their heads firmly stuck in some discarded receptacle are commonplace indicating that if a cat has this information available, it does not always make best use of it.
Marine mammals
Pinnipeds have well-developed tactile senses. Their mystacial vibrissae have ten times the innervation of terrestrial mammals, allowing them to effectively detect vibrations in the water. These vibrations are generated, for example, when a fish swims through water. Detecting vibrations is useful when the animals are foraging and may add to or even replace vision, particularly in darkness.
Harbor seals have been observed following varying paths of other organisms that swam ahead several minutes before, similar to a dog following a scent trail, and even to discriminate the species and the size of the fish responsible for the trail. Blind ringed seals have even been observed successfully hunting on their own in Lake Saimaa, likely relying on their vibrissae to gain sensory information and catch prey. Unlike terrestrial mammals, such as rodents, pinnipeds do not move their vibrissae over an object when examining it but instead extend their moveable whiskers and keep them in the same position. By holding their vibrissae steady, pinnipeds are able to maximize their detection ability. The vibrissae of seals are undulated and wavy while sea lion and walrus vibrissae are smooth. Research is ongoing to determine the function, if any, of these shapes on detection ability. The vibrissa's angle relative to the flow, and not the fiber shape, however, seems to be the most important factor.
Most cetaceans have whiskers at birth but they are typically lost during maturation. The follicles and any vestigial hair sometimes function as touch or electrical sense organs.
Lines of research
Neuroscience
A large part of the brain of whisker-specialist mammals is involved in the processing of nerve impulses from vibrissae, a fact that presumably corresponds to the important position the sense occupies for the animal. Information from the vibrissae arrives in the brain via the trigeminal nerve and is delivered first into the trigeminal sensory complex of brainstem. From there, the most studied pathways are those leading up through parts of thalamus and into barrel cortex, though other major pathways through the superior colliculus in midbrain (a major visual structure in visual animals) and the cerebellum, to name but a couple, are increasingly coming under scrutiny. Neuroscientists, and other researchers, studying sensory systems favour the whisker system for a number of reasons (see Barrel cortex), not least the simple fact that laboratory rats and mice are whisker, rather than visual, specialists.
Evolutionary biology
The presence of mystacial vibrissae in distinct lineages (Rodentia, Afrotheria, marsupials) with remarkable conservation of operation suggests that they may be an old feature present in a common ancestor of all therian mammals. Indeed, some humans even still develop vestigial vibrissal muscles in the upper lip, consistent with the hypothesis that previous members of the human lineage had mystacial vibrissae. Thus, it is possible that the development of the whisker sensory system played an important role in mammalian development, more generally.
Artificial whiskers
Researchers have begun to build artificial whiskers of a variety of types, both to help them understand how biological whiskers work and as a tactile sense for robots. These efforts range from the abstract, through feature-specific models, to attempts to reproduce complete whiskered animals in robot form (ScratchBot and ShrewBot, both robots by Bristol Robotics Laboratory).
In non-mammals
A range of non-mammals possess structures which resemble or function similarly to mammalian whiskers.
In birds
Some birds possess specialized hair-like feathers called rictal bristles around the base of the beak which are sometimes referred to as whiskers.
The whiskered auklet (Aethia pygmaea) has striking, stiff white feathers protruding from above and below the eyes of the otherwise slate-grey bird, and a dark plume which swoops forward from the top of its head. Whiskered auklets sent through a maze of tunnels with their feathers taped back bumped their heads more than twice as often as they did when their feathers were free, indicating they use their feathers in a similar way to cats.
Other birds that have obvious "whiskers" are kiwis, flycatchers, swallows, nightjars, whip-poor-wills, the kākāpō and the long-whiskered owlet (Xenoglaux loweryi).
In fish
Some fish have slender, pendulous tactile organs near the mouth. These are often referred to as "whiskers", although they are more correctly termed barbels. Fish that have barbels include the catfish, carp, goatfish, hagfish, sturgeon, zebrafish and some species of shark.
The Pimelodidae are a family of catfishes (order Siluriformes) commonly known as the long-whiskered catfishes.
In pterosaurs
Anurognathid pterosaurs had a rugose (wrinkled) jaw texture that has been interpreted as the attachment sites for vibrissae, though actual vibrissae have not been recorded. More recently, a specific type of feathers has been found around anurognathid mouths.
Gallery
| Biology and health sciences | Integumentary system | Biology |
251087 | https://en.wikipedia.org/wiki/Gold%20%28color%29 | Gold (color) | Gold, also called golden, is a color tone resembling the gold chemical element.
The web color gold is sometimes referred to as golden to distinguish it from the color metallic gold. The use of gold as a color term in traditional usage is more often applied to the color "metallic gold" (shown below).
The first recorded use of golden as a color name in English was in 1300 to refer to the element gold. The word gold as a color name was first used in 1400 and in 1423 to refer to blond hair.
Metallic gold, such as in paint, is often called goldtone or gold tone, or gold ground when describing a solid gold background. In heraldry, the French word or is used. In model building, the color gold is different from brass. A shiny or metallic silvertone object can be painted with transparent yellow to obtain goldtone, something often done with Christmas decorations.
Metallic gold
Gold (metallic gold)
At right is displayed a representation of the color metallic gold (the color traditionally known as gold) which is a simulation of the color of the actual metallic element gold itself—gold shade.
The source of this color is the ISCC-NBS Dictionary of Color Names (1955), a color dictionary used by stamp collectors to identify the colors of stamps—See color sample of the color Gold (Color Sample Gold (T) #84) displayed on indicated web page:
Web color gold vs. metallic gold
The American Heritage Dictionary defines the color metallic gold as "A light olive-brown to dark yellow, or a moderate, strong to vivid yellow."
Of course, the visual sensation usually associated with the metal gold is its metallic shine. This cannot be reproduced by a simple solid color, because the shiny effect is due to the material's reflective brightness varying with the surface's angle to the light source.
This is why, in art, a metallic paint that glitters in an approximation of real gold would be used; a solid color like that of the cell displayed in the adjacent box does not aesthetically "read" as gold. Especially in sacral art in Christian churches, real gold (as gold leaf) was used for rendering gold in paintings, e.g. for the halo of saints. Gold can also be woven into sheets of silk to give an East Asian traditional look.
More recent art styles, e.g. Art Nouveau, also made use of a metallic, shining gold; however, the metallic finish of such paints was added using fine aluminum powder and pigment rather than actual gold.
Shades
Old gold
Old gold is a dark yellow, which varies from heavy olive or olive brown to deep or strong yellow. The widely accepted color old gold is on the darker rather than the lighter side of this range.
The first recorded use of old gold as a color name in English was in the early 19th century (exact year uncertain). The official colors of Alpha Phi Alpha fraternity, founded in 1906 are black and old gold. The Delta Sigma Pi fraternity, founded in 1907, official colors are designated royal purple and old gold and Pi Kappa Alpha fraternity's colors are garnet and old gold.
Maroon and old gold are the colors of Texas State University's intercollegiate sports teams. Old Gold and black are the team colors of Purdue University Boilermakers intercollegiate sports teams. The Georgia Tech Yellow Jackets formerly wore white and old gold (now called Tech Gold). The Wake Forest Demon Deacons, UCF Knights, and Vanderbilt Commodores wear old gold and black. The UAB Blazers team colors are Forest Green and Old Gold. The New Orleans Saints list their official team colors as black, old gold and white.
Golden yellow
Golden yellow is the color halfway between amber and yellow on the RGB color wheel. It is a color that is 87.5% yellow and 12.5% red.
The first recorded use of golden yellow as a color name in English was in the year 1597.
Golden Yellow is one of the colors of the United States Air Force, along with Ultramarine Blue.
Golden poppy
Golden poppy is a tone of gold that is the color of the California poppy—the official state flower of California—the Golden State.
The first recorded use of golden poppy as a color name in English was in 1927.
Arizona State University (ASU) Gold
Gold is the oldest color associated with Arizona State University and dates back to 1896 when the school was named the Tempe Normal School. Gold signifies the "golden promise" of ASU. Gold also signifies the sunshine Arizona is famous for, including the power of the sun and its influence on the climate and the economy. The student section, known as The Inferno, wears gold on game days.
University of Southern California (USC) Gold
The official colors of the University of Southern California are Pantone 201C and Pantone 123C. These colors, designated as USC Cardinal and USC Gold, were adopted in 1895 by Rev. George W. White, USC's third president, and are equal in importance in identifying the USC Trojans.
California (Berkeley) Gold
California Gold is one of the official colors of the University of California, Berkeley, as identified in their graphic style guide for use in on-screen representations of the gold color in the university's seal. For print media, the guide recommends to, "[u]se Pantone 7750 metallic or Pantone 123 yellow and 282 blue". The color is one of two most used by Berkeley, the other being Berkeley Blue; these, together, are the original colors of the University of California system, of which variations of blue and gold can be found in each campus' school colors.
Cal Poly Pomona gold
Cal Poly Pomona gold was one of the two official colors of California State Polytechnic University, Pomona (Cal Poly Pomona). The official university colors were green (PMS 349) and gold (PMS 131). Cal Poly Pomona's Office of Public Affairs created the colors for web development and has technical guidelines, copyright and privacy protection; as well as logos and images that developers are asked to follow in the university's Guidelines for using official Cal Poly Pomona logos. If web developers are using gold on a university website, they were encouraged to use Cal Poly Pomona gold. Cal Poly Pomona has adopted a new brand color palette including a different gold color: #FFB500. The logo of the Cal Poly Pomona's athletic teams, the Cal Poly Pomona Broncos, has changed in 2014 to reflect the new gold color, but is currently using #FFB718.
UCLA Gold
The color was approved by the University of California, Los Angeles (UCLA) Chancellor in October 2013. This is a shade of gold identified by the university for use in their printed publications.
MU Gold
MU Gold is used by the University of Missouri as the official school color along with black. Mizzou Identity Standards designated the color for web development as well as logos and images that developers are asked to follow in the university's Guidelines for using official Mizzou logos.
Pale gold
The color pale gold is displayed at right.
This has been the color called gold in Crayola crayons since 1903.
Pale gold is one of the Lithuanian basketball club Lietkabelis Panevėžys primary colors.
Gold medal
Gold medal is a tone of metallic gold included in Metallic FX crayons. It was introduced in 2019.
Sunglow
The color sunglow is displayed at right.
This is a Crayola crayon color formulated in 1990.
Harvest gold
The color harvest gold is displayed at right.
This color was originally called harvest in the 1920s.
The first recorded use of harvest as a color name in English was in 1923.
Harvest gold was a common color for metal surfaces (including automobiles and household appliances), as was the color avocado, during the whole decade of the 1970s. They were both also popular colors for shag carpets. Both colors (as well as shag carpets) went out of style by the early 1980s.
Goldenrod
Displayed at right is the web color goldenrod.
The color goldenrod is a representation of the color of some of the deeper gold colored goldenrod flowers.
The first recorded use of goldenrod as a color name in English was in 1915.
Vegas gold
Displayed at right is the color Vegas gold.
Vegas gold, rendered within narrow limits, is associated with the glamorous casinos and hotels of the Las Vegas Strip, United States.
Vegas gold is one of the official athletic colors for the Notre Dame Fighting Irish, Boston College Eagles, Colorado Buffaloes, South Florida Bulls, St. Vincent–St. Mary High School, Vanderbilt Commodores, the United States Naval Academy Midshipmen, and Western Carolina University Catamounts. It is one of the official colors of the NHL's Vegas Golden Knights, and was the type of gold the Pittsburgh Penguins used on their uniforms until they reverted to "Pittsburgh gold", the shade traditionally associated with the city.
Satin sheen gold
At right is displayed the color satin sheen gold. This is the name of the color of the Starfleet command personnel uniform worn by Captain Kirk of the USS Enterprise in the TV show and movies Star Trek.
Golden brown
The first recorded use of golden brown as a color name in English was in the year 1891. Golden brown is commonly referenced in recipes as the desired color of properly baked and fried foods.
Candlelight
Candlelight is a brilliant gold color.
Golden in nature
Protista
The golden algae or chrysophytes are a large group of heterokont algae, found mostly in freshwater.
Plants
Golden bamboo (Phyllostachys aurea) is a bamboo species.
The golden poppy and goldenrod are popular flowers to cultivate in horticulture.
The Yukon Gold potato is a variety of potato recognizable through its smooth eyes and golden interior.
Animals
The golden bamboo lemur (Hapalemur aureus) is a medium-sized bamboo lemur endemic to southeastern Madagascar.
The golden eagle is a Northern Hemisphere bird of prey.
The goldfish was one of the earliest fish to be domesticated, and is still one of the most commonly kept aquarium fish and water garden fish.
The golden jackal is a medium-sized wild canine found in Asia and East Europe
The golden retriever is a medium-sized breed of dog that is one of the most popular companion animals.
The golden toad was an amphibian that used to live in Costa Rica that is now extinct.
In culture
Business
In advertising for the Union Pacific Railroad in the 1950s, the southwestern states of the United States served by the Union Pacific were collectively called The Golden Empire because the railroad's diesel engines were and are colored golden, red, and black. Ads with maps showing the Union Pacific's Golden Empire colored golden were placed in many popular mass-circulation magazines.
Food
Golden rice is a variety of rice produced through genetic engineering to biosynthesize the precursors of beta-carotene (pro-vitamin A) in the edible parts of rice.
Golden Oreos are composed of vanilla instead of chocolate cookies with a vanilla cream filling.
Golden raisins are dried grapes that have been treated with sulfur dioxide and flame-dried.
Gemstones
South Sea Pearls, which have historically been cultured in the Indian and Pacific Oceans, in the countries of Myanmar, Indonesia, the Philippines, and Northern Australia but mostly attributed to the former thalassocratic Sultanate of Sulu have a gold colored variety from the Pinctada maxima Pearl oyster. This golden pearl is the national gemstone of the Philippines. This can now be manufactured in the laboratory at a much lower cost.
Interior design
The Chrysotriklinos (golden reception hall) throne room of the Byzantine Emperor in the Great Palace of Constantinople from its construction, in the late 6th century, until the 10th century
The Queen's Bedchamber in the Grand appartement de la reine in the Versailles Palace is decorated in the color gold. This room was where Marie Antoinette, wife of King Louis XVI of France, slept.
Golden is a warm color that can both provide not only a bright and cheerful feeling but also a somber, traditional, and religious aura. Golden tends to go well with earth colors, but it can also enrich a palette of red or burgundy.
Literature
Blonde hair in women (or sometimes men) is sometimes referred to poetically as golden.
Music
"Golden Brown" is a song by The Stranglers
Parapsychology
Psychics who claim to be able to observe the aura with their third eye report that great spiritual teachers usually have golden auras.
People who have gold auras are said to be those whose pure intellect is applied to abstract philosophy and mathematics.
Politics
The Gold Shirts were a Mexican fascist party in the 1930s.
Gold is often used as an official color by laissez-faire or libertarian political parties, such as the United States Libertarian Party, as well as ideologies such as voluntaryism and anarcho-capitalism, due to their frequent support for Austrian School economics and the gold standard.
Religion
The color golden is associated with Buddhism:
Statues of Buddha are usually painted metallic gold, are made of the metal gold, or have gold plating.
Theravada Buddhist monks wear saffron robes, a color close to golden.
The Shwedagon Pagoda, in Yangon, Myanmar is plated in solid gold.
The Golden Pavilion is a notable Buddhist temple in Kyoto, Kinki, Japan.
The Secret of the Golden Flower is an important religious text in Daoism.
The Golden Temple in Amritsar, Punjab, India, is the holiest site of the Sikh religion.
The Golden Mosque in Samarra, Iraq, a Shiite Muslim holy site constructed in 944.
In Paganism, it is used for inner strength, self-realization, understanding, and intuition. It has a masculine energy and brings fortune and luck.
Sports
In association football, the Wolverhampton Wanderers traditionally uses "old gold" as its primary color, though the shade of the color is not quite metallic.
In Major League Baseball, the Milwaukee Brewers, Oakland Athletics, and Pittsburgh Pirates use gold as one of their team's primary colors. The three teams utilize a more yellowish shade of the color (from 1994 to 2019, the Brewers used a darker gold before reverting to the yellowish hue in 2020). A Gold Glove Award is given to the best player at each fielding position in each major league.
The Kansas City Royals use gold as an accent color.
In the NBA, the Boston Celtics use "old gold" as an accent color. The Cleveland Cavaliers, Los Angeles Lakers, Indiana Pacers, Denver Nuggets, Golden State Warriors, Memphis Grizzlies, New Orleans Pelicans, Utah Jazz and Oklahoma City Thunder also use gold as an accent color, although the shade is mostly yellowish in appearance.
Fans of the National Football League will note the Los Angeles Rams, Pittsburgh Steelers, Washington Commanders, Minnesota Vikings, Green Bay Packers, and Kansas City Chiefs as having gold as a color. The gold they use, however, is a distinctly more yellow color (akin to the non-metallic web color version) than the traditional "old gold" used by the New Orleans Saints and San Francisco 49ers. In both cases, the color is referred to as "gold", with the yellow shade sometimes referred to as "athletic gold" when distinguishing it from the metallic shades.
The Rams used a darker gold which was called "New Century Gold" from 2000 to 2015 when they played in St. Louis.
In the NHL, eleven teams currently use a form of gold in their color schemes. The Boston Bruins are the oldest team to do so, and have always used a yellow "athletic" gold; they have even worn several gold sweaters throughout the years. Other teams using athletic gold include the St. Louis Blues, Buffalo Sabres (excluding the 1996–2006 period), Calgary Flames, Florida Panthers (who use a darker, more metallic shade in their logo) and Nashville Predators. Teams that have switched from athletic gold to a metallic shade include the Minnesota North Stars (now the Dallas Stars), Pittsburgh Penguins, and Anaheim Ducks. The Ottawa Senators and Minnesota Wild have always used metallic gold. In addition, the Los Angeles Kings and Vancouver Canucks have used versions of athletic gold in the past. The Kings, Canucks, and Penguins have also worn gold sweaters color in the past. For the 2014–15 season, the Penguins revived the black uniform the team wore during its first two championship seasons in 1991 and 1992, with the team's old shade of gold as an alternate uniform. When the team first switched from shades of blue to black and gold in 1980, the color scheme was similar to that of the Bruins, who protested the Penguins' new uniforms. The protest failed—largely due to Pittsburgh having a tradition of sports teams wearing gold and black—and the Penguins wore the shade of gold, now dubbed "Pittsburgh gold", until 2002 and once again in 2014.
In college sports, the U.S. Military Academy and the U.S. Naval Academy use gold as a primary color.
Vexillology
Argentina, Brazil, Bolivia, Bosnia and Herzegovina, Belgium, Bhutan, China, Colombia, Egypt, Ecuador, Germany, Ghana, Malaysia, the Philippines Spain, Sri Lanka, Venezuela, Vietnam are examples of modern nations that use the color golden in their national flags.
The Holy Roman Empire, which existed from 800 to 1806, had a golden flag with a black double-headed Imperial Eagle on the field, the origin of the use of the color in the German and Belgian flags.
The Byzantine Empire from 1261 until its collapse in 1453 had a flag that had a black double-headed eagle on a field of golden. This flag is still used today as the flag of the Mount Athos autonomous region in Greece.
The Flag of the Hispanic People (Bandera de la Raza) is an ethnic flag that is golden and purpure (purple) on a white field. It is also used as the flag of Hispanic America. (This flag is sometimes also called the Flag of the Americas when used on a non-ethnic basis to symbolically represent the combined geographical area of North America and South America together.)
Chemical compounds
In addition to elemental gold, a number of compounds or alloys have a reflective gold hue:
Several brasses, specifically those rich (65%+ wt.) in copper. Prince's metal is a brass-based gold simulant.
Titanium nitride
Zirconium nitride
| Physical sciences | Colors | Physics |
251147 | https://en.wikipedia.org/wiki/Isatis%20tinctoria | Isatis tinctoria | Isatis tinctoria, also called woad (), dyer's woad, dyer's-weed, or glastum, is a flowering plant in the family Brassicaceae (the mustard family) with a documented history of use as a blue dye and medicinal plant.
Its genus name, Isatis, derives from the ancient Greek word for the plant, . It is occasionally known as Asp of Jerusalem. Woad is also the name of a blue dye produced from the leaves of the plant. Woad is native to the steppe and desert zones of the Caucasus, Central Asia to Eastern Siberia and Western Asia but is now also found in South-Eastern and Central Europe and western North America.
Since ancient times, woad was an important source of blue dye and was cultivated throughout Europe, especially in Western and Southern Europe. In medieval times, there were important woad-growing regions in England, Germany and France. Towns such as Toulouse became prosperous from the woad trade.
Woad was eventually replaced by the more colourfast Indigofera tinctoria and, in the early 20th century, both woad and Indigofera tinctoria were replaced by synthetic blue dyes. Woad has been used medicinally for centuries. The double use of woad is seen in its name: the term Isatis is linked to its ancient use to treat wounds; the term tinctoria references its use as a dye. There has also been some revival of the use of woad for craft purposes.
History of woad cultivation
Ancient use
The first archaeological finds of woad seeds date to the Neolithic period. The seeds have been found in the cave of l'Audoste, Bouches-du-Rhône, France. Impressions of seeds of Färberwaid (Isatis tinctoria L.) or German indigo, of the plant family Brassicaceae, have been found on pottery in the Iron Age settlement of Heuneburg, Germany. Seed and pod fragments have also been found in an Iron Age pit at Dragonby, North Lincolnshire, United Kingdom. The Hallstatt burials of the Hochdorf Chieftain's Grave and Hohmichele contained textiles dyed with woad.
Melo and Rondão write that woad was known "as far back as the time of the ancient Egyptians, who used it to dye the cloth wrappings applied for the mummies." Skelton states that one of the early dyes discovered by the ancient Egyptians was "blue woad (Isatis tinctoria)." Lucas writes, "What has been assumed to have been Indian Indigo on ancient Egyptian fabrics may have been woad." Hall states that the ancient Egyptians created their blue dye "by using indigotin, otherwise known as woad."
A dye known as in Aramaic, is mentioned in the Babylonian Talmud.
Celtic blue is a shade of blue, also known as in Welsh, or in both the Irish language and in Scottish Gaelic. Julius Caesar reported (in ) that the Britanni used to colour their bodies blue with , a word that means primarily , but also the domestic name for the woad (Isatis tinctoria), besides the Gaulish loanword (from Proto-Celtic ). The connection seems to be that both glass and the woad are "water-like" ( is from Proto-Indo-European , ). In terms of usage, the Latin is more often used to refer to glass rather than woad. The use of the word for the woad might also be understood as "coloured like glass", applied to the plant and the dye made from it.
Gillian Carr conducted experiments using indigo pigment derived from woad mixed with different binders to make body paint. The resulting paints yielded colours from "grey-blue, through intense midnight blue, to black". People with modern experiences with woad as a tattoo pigment have claimed that it does not work well, and is actually caustic and causes scarring when put into the skin.
It has also been claimed that Caesar was referring to some form of copper- or iron-based pigment. Analysis done on the Lindow Man did return evidence of copper. The same study also noted that the earliest definite reference to the woad plant in the British Isles dates to a seed impression on an Anglo-Saxon pot. The authors theorize that vitrum could have actually referred to copper(II) sulfate's naturally occurring variant chalcanthite or to the mineral azurite. A later study concluded the amount was "not of sufficient magnitude to provide convincing evidence that the copper was deliberately applied as paint".
Woad was an important dyeing agent in much of Europe and parts of England during the medieval period. However, dye traders began to import indigo during the sixteenth and seventeenth centuries, which threatened to replace locally grown woad as the primary blue dye. The translation of vitrum as woad may date to this period.
Medieval period onwards
Woad was one of the three staples of the European dyeing industry, along with weld (yellow) and madder (red). Chaucer mentions their use by the dyer ("litestere") in his poem The Former Age:
No mader, welde, or wood no litestere
Ne knew; the flees was of his former hewe;
The three colours can be seen together in tapestries such as The Hunt of the Unicorn (1495–1505), though typically it is the dark blue of the woad that has lasted best. Medieval uses of the dye were not limited to textiles. For example, the illustrator of the Lindisfarne Gospels () used a woad-based pigment for blue paint. As does the late 13th century North Italian manual on book illumination Liber colorum secundum magistrum Bernardum describe its usage.
In Viking Age levels at archaeological digs at York, a dye shop with remains of both woad and madder have been excavated and dated to the 10th century. In medieval times, centres of woad cultivation lay in Lincolnshire and Somerset in England, Jülich and the Erfurt area in Thuringia in Germany, Piedmont and Tuscany in Italy, and Gascogne, Normandy, the Somme Basin (from Amiens to Saint-Quentin), Brittany and, above all, Languedoc in France. This last region, in the triangle created by Toulouse, Albi and Carcassonne, known as the Lauragais, was for a long time the biggest producer of woad, or pastel, as it was locally known. One writer commented that "woad […] hath made that country the happiest and richest in Europe."
The prosperous woad merchants of Toulouse displayed their affluence in splendid mansions, many of which still stand, as the Hôtel de Bernuy and the Hôtel d'Assézat. One merchant, Jean de Bernuy, a Spanish Jew who had fled the Spanish Inquisition, was credit-worthy enough to be the main guarantor of the ransomed King Francis I after his capture at the Battle of Pavia by Charles V of Spain. Much of the woad produced here was used for the cloth industry in southern France, but it was also exported via Bayonne, Narbonne and Bordeaux to Flanders, the Low Countries, Italy, and above all Britain and Spain.
After cropping the woad eddish could be let out for grazing sheep.
The woad produced in Lincolnshire and Cambridgeshire in the 19th century was shipped out from the Port of Wisbech, Spalding and Boston, both the last to northern mills and the USA. The last portable woad mill was at Parson Drove, Cambridgeshire, Wisbech & Fenland Museum has a woad mill model, photos and other items used in woad production.
A major market for woad was at Görlitz in Lausitz. The citizens of the five Thuringian Färberwaid (dye woad) towns of Erfurt, Gotha, Tennstedt, Arnstadt and Langensalza had their own charters. In Erfurt, the woad-traders gave the funds to found the University of Erfurt. Traditional fabric is still printed with woad in Thuringia, Saxony and Lusatia today: it is known as Blaudruck (literally, "blue print(ing)").
In the Marche region, the cultivation of the plant was an important resource for the Duchy of Urbino in Italy. To fully understand the importance of the ford industry in the State of Urbino, it is enough to read the comprehensive Chapters of the art of wool in 1555, which dictated prescriptions regarding the cultivation and trade of woad, whether in loaves or macerated (powdered). Testifying to the importance that this crop had in the economy in addition to the archival documents was the identification of a hundred millstones surveyed by Delio Bischi in the province of Pesaro and Urbino, the original use of which had become completely unknown as their memory had been lost.
Use as Chinese medicine
The woad plant's roots are used in Traditional Chinese medicine to make a medicine known as banlangen ( ) that purports to have antiviral properties. Banlangen is used as an herbal medicinal tea in China for colds and tonsillar ailments. Used as a tea, it has a brownish appearance and (unlike most Chinese medicines) is mildly sweet in taste.
Woad and indigo
The dye chemical extracted from woad is indigo, the same dye extracted from "true indigo", Indigofera tinctoria, but in a lower concentration. Following the Portuguese discovery of the sea route to India by the navigator Vasco da Gama in 1498, great amounts of indigo were imported from Asia. Laws were passed in some parts of Europe to protect the woad industry from the competition of the indigo trade. It was proclaimed that indigo caused yarns to rot. This prohibition was repeated in 1594 and again in 1603. In France, Henry IV, in an edict of 1609, forbade under pain of death the use of "the false and pernicious Indian drug".
With the development of a chemical process to synthesize the pigment, both the woad and natural indigo industries collapsed in the first years of the 20th century. The last commercial harvest of woad until recent times occurred in 1932, in Lincolnshire, Britain. Small amounts of woad are now grown in the UK and France to supply craft dyers. The classic book about woad is The Woad Plant and its Dye by J. B. Hurry, Oxford University Press of 1930, which contains an extensive bibliography.
A method for producing blue dye from woad is described in The History of Woad and the Medieval Woad Vat (1998) .
Woad is biodegradable and safe in the environment. In Germany, there have been attempts to use it to protect wood against decay without applying dangerous chemicals. Production of woad is increasing in the UK for use in inks, particularly for inkjet printers, and dyes.
Invasive and noxious weed
In certain locations, the plant is classified as a non-native and invasive weed. It is listed as a noxious weed by the agriculture departments of several states in the western United States: Arizona, California, Colorado, Idaho, Montana, Nevada, New Mexico, Oregon, Utah, Washington, and Wyoming. In Montana, it has been the target of an extensive, and largely successful, eradication attempt.
| Biology and health sciences | Brassicales | Plants |
251399 | https://en.wikipedia.org/wiki/Observable%20universe | Observable universe | The observable universe is a spherical region of the universe consisting of all matter that can be observed from Earth; the electromagnetic radiation from these objects has had time to reach the Solar System and Earth since the beginning of the cosmological expansion. Assuming the universe is isotropic, the distance to the edge of the observable universe is roughly the same in every direction. That is, the observable universe is a spherical region centered on the observer. Every location in the universe has its own observable universe, which may or may not overlap with the one centered on Earth.
The word observable in this sense does not refer to the capability of modern technology to detect light or other information from an object, or whether there is anything to be detected. It refers to the physical limit created by the speed of light itself. No signal can travel faster than light, hence there is a maximum distance, called the particle horizon, beyond which nothing can be detected, as the signals could not have reached us yet. Sometimes astrophysicists distinguish between the observable universe and the visible universe. The former includes signals since the end of the inflationary epoch, while the latter includes only signals emitted since recombination.
According to calculations, the current comoving distance to particles from which the cosmic microwave background radiation (CMBR) was emitted, which represents the radius of the visible universe, is about 14.0 billion parsecs (about 45.7 billion light-years). The comoving distance to the edge of the observable universe is about 14.3 billion parsecs (about 46.6 billion light-years), about 2% larger. The radius of the observable universe is therefore estimated to be about 46.5 billion light-years. Using the critical density and the diameter of the observable universe, the total mass of ordinary matter in the universe can be calculated to be about . In November 2018, astronomers reported that extragalactic background light (EBL) amounted to photons.
As the universe's expansion is accelerating, all currently observable objects, outside the local supercluster, will eventually appear to freeze in time, while emitting progressively redder and fainter light. For instance, objects with the current redshift z from 5 to 10 will only be observable up to an age of 4–6 billion years. In addition, light emitted by objects currently situated beyond a certain comoving distance (currently about ) will never reach Earth.
Overview
The universe's size is unknown, and it may be infinite in extent. Some parts of the universe are too far away for the light emitted since the Big Bang to have had enough time to reach Earth or space-based instruments, and therefore lie outside the observable universe. In the future, light from distant galaxies will have had more time to travel, so one might expect that additional regions will become observable. Regions distant from observers (such as us) are expanding away faster than the speed of light, at rates estimated by Hubble's law. The expansion rate appears to be accelerating, which dark energy was proposed to explain.
Assuming dark energy remains constant (an unchanging cosmological constant) so that the expansion rate of the universe continues to accelerate, there is a "future visibility limit" beyond which objects will never enter the observable universe at any time in the future because light emitted by objects outside that limit could never reach the Earth. Note that, because the Hubble parameter is decreasing with time, there can be cases where a galaxy that is receding from Earth only slightly faster than light emits a signal that eventually reaches Earth. This future visibility limit is calculated at a comoving distance of 19 billion parsecs (62 billion light-years), assuming the universe will keep expanding forever, which implies the number of galaxies that can ever be theoretically observed in the infinite future is only larger than the number currently observable by a factor of 2.36 (ignoring redshift effects).
In principle, more galaxies will become observable in the future; in practice, an increasing number of galaxies will become extremely redshifted due to ongoing expansion, so much so that they will seem to disappear from view and become invisible. A galaxy at a given comoving distance is defined to lie within the "observable universe" if we can receive signals emitted by the galaxy at any age in its history, say, a signal sent from the galaxy only 500 million years after the Big Bang. Because of the universe's expansion, there may be some later age at which a signal sent from the same galaxy can never reach the Earth at any point in the infinite future, so, for example, we might never see what the galaxy looked like 10 billion years after the Big Bang, even though it remains at the same comoving distance less than that of the observable universe.
This can be used to define a type of cosmic event horizon whose distance from the Earth changes over time. For example, the current distance to this horizon is about 16 billion light-years, meaning that a signal from an event happening at present can eventually reach the Earth if the event is less than 16 billion light-years away, but the signal will never reach the Earth if the event is further away.
The space before this cosmic event horizon can be called "reachable universe", that is all galaxies closer than that could be reached if we left for them today, at the speed of light; all galaxies beyond that are unreachable. Simple observation will show the future visibility limit (62 billion light-years) is exactly equal to the reachable limit (16 billion light-years) added to the current visibility limit (46 billion light-years).
"The universe" versus "the observable universe"
Both popular and professional research articles in cosmology often use the term "universe" to mean "observable universe". This can be justified on the grounds that we can never know anything by direct observation about any part of the universe that is causally disconnected from the Earth, although many credible theories require a total universe much larger than the observable universe. No evidence exists to suggest that the boundary of the observable universe constitutes a boundary on the universe as a whole, nor do any of the mainstream cosmological models propose that the universe has any physical boundary in the first place. However, some models propose it could be finite but unbounded, like a higher-dimensional analogue of the 2D surface of a sphere that is finite in area but has no edge.
It is plausible that the galaxies within the observable universe represent only a minuscule fraction of the galaxies in the universe. According to the theory of cosmic inflation initially introduced by Alan Guth and D. Kazanas, if it is assumed that inflation began about 10−37 seconds after the Big Bang and that the pre-inflation size of the universe was approximately equal to the speed of light times its age, that would suggest that at present the entire universe's size is at least light-years—at least times the radius of the observable universe.
If the universe is finite but unbounded, it is also possible that the universe is smaller than the observable universe. In this case, what we take to be very distant galaxies may actually be duplicate images of nearby galaxies, formed by light that has circumnavigated the universe. It is difficult to test this hypothesis experimentally because different images of a galaxy would show different eras in its history, and consequently might appear quite different. Bielewicz et al. claim to establish a lower bound of 27.9 gigaparsecs (91 billion light-years) on the diameter of the last scattering surface. This value is based on matching-circle analysis of the WMAP 7-year data. This approach has been disputed.
Size
The comoving distance from Earth to the edge of the observable universe is about 14.26 gigaparsecs (46.5 billion light-years or ) in any direction. The observable universe is thus a sphere with a diameter of about 28.5 gigaparsecs (93 billion light-years or ). Assuming that space is roughly flat (in the sense of being a Euclidean space), this size corresponds to a comoving volume of about ( or ).
These are distances now (in cosmological time), not distances at the time the light was emitted. For example, the cosmic microwave background radiation that we see right now was emitted at the time of photon decoupling, estimated to have occurred about years after the Big Bang, which occurred around 13.8 billion years ago. This radiation was emitted by matter that has, in the intervening time, mostly condensed into galaxies, and those galaxies are now calculated to be about 46 billion light-years from Earth. To estimate the distance to that matter at the time the light was emitted, we may first note that according to the Friedmann–Lemaître–Robertson–Walker metric, which is used to model the expanding universe, if we receive light with a redshift of z, then the scale factor at the time the light was originally emitted is given by
.
WMAP nine-year results combined with other measurements give the redshift of photon decoupling as z = , which implies that the scale factor at the time of photon decoupling would be . So if the matter that originally emitted the oldest CMBR photons has a present distance of 46 billion light-years, then the distance would have been only about 42 million light-years at the time of decoupling.
The light-travel distance to the edge of the observable universe is the age of the universe times the speed of light, 13.8 billion light years. This is the distance that a photon emitted shortly after the Big Bang, such as one from the cosmic microwave background, has traveled to reach observers on Earth. Because spacetime is curved, corresponding to the expansion of space, this distance does not correspond to the true distance at any moment in time.
Matter and mass
Number of galaxies and stars
The observable universe contains as many as an estimated 2 trillion galaxies and, overall, as many as an estimated 1024 stars – more stars (and, potentially, Earth-like planets) than all the grains of beach sand on planet Earth. Other estimates are in the hundreds of billions rather than trillions. The estimated total number of stars in an inflationary universe (observed and unobserved) is 10100.
Matter content—number of atoms
Assuming the mass of ordinary matter is about as discussed above, and assuming all atoms are hydrogen atoms (which are about 74% of all atoms in the Milky Way by mass), the estimated total number of atoms in the observable universe is obtained by dividing the mass of ordinary matter by the mass of a hydrogen atom. The result is approximately 1080 hydrogen atoms, also known as the Eddington number.
Mass of ordinary matter
The mass of the observable universe is often quoted as 1053 kg. In this context, mass refers to ordinary (baryonic) matter and includes the interstellar medium (ISM) and the intergalactic medium (IGM). However, it excludes dark matter and dark energy. This quoted value for the mass of ordinary matter in the universe can be estimated based on critical density. The calculations are for the observable universe only as the volume of the whole is unknown and may be infinite.
Estimates based on critical density
Critical density is the energy density for which the universe is flat. If there is no dark energy, it is also the density for which the expansion of the universe is poised between continued expansion and collapse. From the Friedmann equations, the value for critical density, is:
where G is the gravitational constant and is the present value of the Hubble constant. The value for H0, as given by the European Space Agency's Planck Telescope, is H0 = 67.15 kilometres per second per megaparsec. This gives a critical density of , or about 5 hydrogen atoms per cubic metre. This density includes four significant types of energy/mass: ordinary matter (4.8%), neutrinos (0.1%), cold dark matter (26.8%), and dark energy (68.3%).
Although neutrinos are Standard Model particles, they are listed separately because they are ultra-relativistic and hence behave like radiation rather than like matter. The density of ordinary matter, as measured by Planck, is 4.8% of the total critical density or . To convert this density to mass we must multiply by volume, a value based on the radius of the "observable universe". Since the universe has been expanding for 13.8 billion years, the comoving distance (radius) is now about 46.6 billion light-years. Thus, volume (πr3) equals and the mass of ordinary matter equals density () times volume () or .
Large-scale structure
Sky surveys and mappings of the various wavelength bands of electromagnetic radiation (in particular 21-cm emission) have yielded much information on the content and character of the universe's structure. The organization of structure appears to follow a hierarchical model with organization up to the scale of superclusters and filaments. Larger than this (at scales between 30 and 200 megaparsecs), there seems to be no continued structure, a phenomenon that has been referred to as the End of Greatness.
Walls, filaments, nodes, and voids
The organization of structure arguably begins at the stellar level, though most cosmologists rarely address astrophysics on that scale. Stars are organized into galaxies, which in turn form galaxy groups, galaxy clusters, superclusters, sheets, walls and filaments, which are separated by immense voids, creating a vast foam-like structure sometimes called the "cosmic web". Prior to 1989, it was commonly assumed that virialized galaxy clusters were the largest structures in existence, and that they were distributed more or less uniformly throughout the universe in every direction. However, since the early 1980s, more and more structures have been discovered. In 1983, Adrian Webster identified the Webster LQG, a large quasar group consisting of 5 quasars. The discovery was the first identification of a large-scale structure, and has expanded the information about the known grouping of matter in the universe.
In 1987, Robert Brent Tully identified the Pisces–Cetus Supercluster Complex, the galaxy filament in which the Milky Way resides. It is about 1 billion light-years across. That same year, an unusually large region with a much lower than average distribution of galaxies was discovered, the Giant Void, which measures 1.3 billion light-years across. Based on redshift survey data, in 1989 Margaret Geller and John Huchra discovered the "Great Wall", a sheet of galaxies more than 500 million light-years long and 200 million light-years wide, but only 15 million light-years thick. The existence of this structure escaped notice for so long because it requires locating the position of galaxies in three dimensions, which involves combining location information about the galaxies with distance information from redshifts.
Two years later, astronomers Roger G. Clowes and Luis E. Campusano discovered the Clowes–Campusano LQG, a large quasar group measuring two billion light-years at its widest point, which was the largest known structure in the universe at the time of its announcement. In April 2003, another large-scale structure was discovered, the Sloan Great Wall. In August 2007, a possible supervoid was detected in the constellation Eridanus. It coincides with the 'CMB cold spot', a cold region in the microwave sky that is highly improbable under the currently favored cosmological model. This supervoid could cause the cold spot, but to do so it would have to be improbably big, possibly a billion light-years across, almost as big as the Giant Void mentioned above.
Another large-scale structure is the SSA22 Protocluster, a collection of galaxies and enormous gas bubbles that measures about 200 million light-years across.
In 2011, a large quasar group was discovered, U1.11, measuring about 2.5 billion light-years across. On January 11, 2013, another large quasar group, the Huge-LQG, was discovered, which was measured to be four billion light-years across, the largest known structure in the universe at that time. In November 2013, astronomers discovered the Hercules–Corona Borealis Great Wall, an even bigger structure twice as large as the former. It was defined by the mapping of gamma-ray bursts.
In 2021, the American Astronomical Society announced the detection of the Giant Arc; a crescent-shaped string of galaxies that span 3.3 billion light years in length, located 9.2 billion light years from Earth in the constellation Boötes from observations captured by the Sloan Digital Sky Survey.
End of Greatness
The End of Greatness is an observational scale discovered at roughly 100 Mpc (roughly 300 million light-years) where the lumpiness seen in the large-scale structure of the universe is homogenized and isotropized in accordance with the cosmological principle. At this scale, no pseudo-random fractalness is apparent.
The superclusters and filaments seen in smaller surveys are randomized to the extent that the smooth distribution of the universe is visually apparent. It was not until the redshift surveys of the 1990s were completed that this scale could accurately be observed.
Observations
Another indicator of large-scale structure is the 'Lyman-alpha forest'. This is a collection of absorption lines that appear in the spectra of light from quasars, which are interpreted as indicating the existence of huge thin sheets of intergalactic (mostly hydrogen) gas. These sheets appear to collapse into filaments, which can feed galaxies as they grow where filaments either cross or are dense. An early direct evidence for this cosmic web of gas was the 2019 detection, by astronomers from the RIKEN Cluster for Pioneering Research in Japan and Durham University in the U.K., of light from the brightest part of this web, surrounding and illuminated by a cluster of forming galaxies, acting as cosmic flashlights for intercluster medium hydrogen fluorescence via Lyman-alpha emissions.
In 2021, an international team, headed by Roland Bacon from the Centre de Recherche Astrophysique de Lyon (France), reported the first observation of diffuse extended Lyman-alpha emission from redshift 3.1 to 4.5 that traced several cosmic web filaments on scales of 2.5−4 cMpc (comoving mega-parsecs), in filamentary environments outside massive structures typical of web nodes.
Some caution is required in describing structures on a cosmic scale because they are often different from how they appear. Gravitational lensing can make an image appear to originate in a different direction from its real source, when foreground objects curve surrounding spacetime (as predicted by general relativity) and deflect passing light rays. Rather usefully, strong gravitational lensing can sometimes magnify distant galaxies, making them easier to detect. Weak lensing by the intervening universe in general also subtly changes the observed large-scale structure.
The large-scale structure of the universe also looks different if only redshift is used to measure distances to galaxies. For example, galaxies behind a galaxy cluster are attracted to it and fall towards it, and so are blueshifted (compared to how they would be if there were no cluster). On the near side, objects are redshifted. Thus, the environment of the cluster looks somewhat pinched if using redshifts to measure distance. The opposite effect is observed on galaxies already within a cluster: the galaxies have some random motion around the cluster center, and when these random motions are converted to redshifts, the cluster appears elongated. This creates a "finger of God"—the illusion of a long chain of galaxies pointed at Earth.
Cosmography of Earth's cosmic neighborhood
At the centre of the Hydra–Centaurus Supercluster, a gravitational anomaly called the Great Attractor affects the motion of galaxies over a region hundreds of millions of light-years across. These galaxies are all redshifted, in accordance with Hubble's law. This indicates that they are receding from us and from each other, but the variations in their redshift are sufficient to reveal the existence of a concentration of mass equivalent to tens of thousands of galaxies.
The Great Attractor, discovered in 1986, lies at a distance of between 150 million and 250 million light-years in the direction of the Hydra and Centaurus constellations. In its vicinity there is a preponderance of large old galaxies, many of which are colliding with their neighbours, or radiating large amounts of radio waves.
In 1987, astronomer R. Brent Tully of the University of Hawaii's Institute of Astronomy identified what he called the Pisces–Cetus Supercluster Complex, a structure one billion light-years long and 150 million light-years across in which, he claimed, the Local Supercluster is embedded.
Most distant objects
The most distant astronomical object identified (as of August of 2024) is a galaxy classified as JADES-GS-z14-0. In 2009, a gamma ray burst, GRB 090423, was found to have a redshift of 8.2, which indicates that the collapsing star that caused it exploded when the universe was only 630 million years old. The burst happened approximately 13 billion years ago, so a distance of about 13 billion light-years was widely quoted in the media, or sometimes a more precise figure of 13.035 billion light-years.
This would be the "light travel distance" (see Distance measures (cosmology)) rather than the "proper distance" used in both Hubble's law and in defining the size of the observable universe. Cosmologist Ned Wright argues against using this measure. The proper distance for a redshift of 8.2 would be about 9.2 Gpc, or about 30 billion light-years.
Horizons
The limit of observability in the universe is set by cosmological horizons which limit—based on various physical constraints—the extent to which information can be obtained about various events in the universe. The most famous horizon is the particle horizon which sets a limit on the precise distance that can be seen due to the finite age of the universe. Additional horizons are associated with the possible future extent of observations, larger than the particle horizon owing to the expansion of space, an "optical horizon" at the surface of last scattering, and associated horizons with the surface of last scattering for neutrinos and gravitational waves.
Gallery
| Physical sciences | Physical cosmology | null |
251478 | https://en.wikipedia.org/wiki/Quadratic%20form | Quadratic form | In mathematics, a quadratic form is a polynomial with terms all of degree two ("form" is another name for a homogeneous polynomial). For example,
is a quadratic form in the variables and . The coefficients usually belong to a fixed field , such as the real or complex numbers, and one speaks of a quadratic form over . Over the reals, a quadratic form is said to be definite if it takes the value zero only when all its variables are simultaneously zero; otherwise it is isotropic.
Quadratic forms occupy a central place in various branches of mathematics, including number theory, linear algebra, group theory (orthogonal groups), differential geometry (the Riemannian metric, the second fundamental form), differential topology (intersection forms of manifolds, especially four-manifolds), Lie theory (the Killing form), and statistics (where the exponent of a zero-mean multivariate normal distribution has the quadratic form )
Quadratic forms are not to be confused with quadratic equations, which have only one variable and may include terms of degree less than two. A quadratic form is specific instance of the more general concept of forms.
Introduction
Quadratic forms are homogeneous quadratic polynomials in variables. In the cases of one, two, and three variables they are called unary, binary, and ternary and have the following explicit form:
where , ..., are the coefficients.
The theory of quadratic forms and methods used in their study depend in a large measure on the nature of the coefficients, which may be real or complex numbers, rational numbers, or integers. In linear algebra, analytic geometry, and in the majority of applications of quadratic forms, the coefficients are real or complex numbers. In the algebraic theory of quadratic forms, the coefficients are elements of a certain field. In the arithmetic theory of quadratic forms, the coefficients belong to a fixed commutative ring, frequently the integers or the -adic integers . Binary quadratic forms have been extensively studied in number theory, in particular, in the theory of quadratic fields, continued fractions, and modular forms. The theory of integral quadratic forms in variables has important applications to algebraic topology.
Using homogeneous coordinates, a non-zero quadratic form in variables defines an -dimensional quadric in the -dimensional projective space. This is a basic construction in projective geometry. In this way one may visualize 3-dimensional real quadratic forms as conic sections.
An example is given by the three-dimensional Euclidean space and the square of the Euclidean norm expressing the distance between a point with coordinates and the origin:
A closely related notion with geometric overtones is a quadratic space, which is a pair , with a vector space over a field , and a quadratic form on V. See below for the definition of a quadratic form on a vector space.
History
The study of quadratic forms, in particular the question of whether a given integer can be the value of a quadratic form over the integers, dates back many centuries. One such case is Fermat's theorem on sums of two squares, which determines when an integer may be expressed in the form , where , are integers. This problem is related to the problem of finding Pythagorean triples, which appeared in the second millennium BCE.
In 628, the Indian mathematician Brahmagupta wrote Brāhmasphuṭasiddhānta, which includes, among many other things, a study of equations of the form . He considered what is now called Pell's equation, , and found a method for its solution. In Europe this problem was studied by Brouncker, Euler and Lagrange.
In 1801 Gauss published Disquisitiones Arithmeticae, a major portion of which was devoted to a complete theory of binary quadratic forms over the integers. Since then, the concept has been generalized, and the connections with quadratic number fields, the modular group, and other areas of mathematics have been further elucidated.
Associated symmetric matrix
Any matrix determines a quadratic form in variables by
where .
Example
Consider the case of quadratic forms in three variables . The matrix has the form
The above formula gives
So, two different matrices define the same quadratic form if and only if they have the same elements on the diagonal and the same values for the sums , and . In particular, the quadratic form is defined by a unique symmetric matrix
This generalizes to any number of variables as follows.
General case
Given a quadratic form over the real numbers, defined by the matrix ,
the matrix
is symmetric, defines the same quadratic form as , and is the unique symmetric matrix that defines .
So, over the real numbers (and, more generally, over a field of characteristic different from two), there is a one-to-one correspondence between quadratic forms and symmetric matrices that determine them.
Real quadratic forms
A fundamental problem is the classification of real quadratic forms under a linear change of variables.
Jacobi proved that, for every real quadratic form, there is an orthogonal diagonalization; that is, an orthogonal change of variables that puts the quadratic form in a "diagonal form"
where the associated symmetric matrix is diagonal. Moreover, the coefficients are determined uniquely up to a permutation.
If the change of variables is given by an invertible matrix that is not necessarily orthogonal, one can suppose that all coefficients are 0, 1, or −1. Sylvester's law of inertia states that the numbers of each 0, 1, and −1 are invariants of the quadratic form, in the sense that any other diagonalization will contain the same number of each. The signature of the quadratic form is the triple , where these components count the number of 0s, number of 1s, and the number of −1s, respectively. Sylvester's law of inertia shows that this is a well-defined quantity attached to the quadratic form.
The case when all have the same sign is especially important: in this case the quadratic form is called positive definite (all 1) or negative definite (all −1). If none of the terms are 0, then the form is called ; this includes positive definite, negative definite, and isotropic quadratic form (a mix of 1 and −1); equivalently, a nondegenerate quadratic form is one whose associated symmetric form is a nondegenerate bilinear form. A real vector space with an indefinite nondegenerate quadratic form of index (denoting 1s and −1s) is often denoted as particularly in the physical theory of spacetime.
The discriminant of a quadratic form, concretely the class of the determinant of a representing matrix in (up to non-zero squares) can also be defined, and for a real quadratic form is a cruder invariant than signature, taking values of only "positive, zero, or negative". Zero corresponds to degenerate, while for a non-degenerate form it is the parity of the number of negative coefficients, .
These results are reformulated in a different way below.
Let be a quadratic form defined on an -dimensional real vector space. Let be the matrix of the quadratic form in a given basis. This means that is a symmetric matrix such that
where x is the column vector of coordinates of in the chosen basis. Under a change of basis, the column is multiplied on the left by an invertible matrix , and the symmetric square matrix is transformed into another symmetric square matrix of the same size according to the formula
Any symmetric matrix can be transformed into a diagonal matrix
by a suitable choice of an orthogonal matrix , and the diagonal entries of are uniquely determined – this is Jacobi's theorem. If is allowed to be any invertible matrix then can be made to have only 0, 1, and −1 on the diagonal, and the number of the entries of each type ( for 0, for 1, and for −1) depends only on . This is one of the formulations of Sylvester's law of inertia and the numbers and are called the positive and negative indices of inertia. Although their definition involved a choice of basis and consideration of the corresponding real symmetric matrix , Sylvester's law of inertia means that they are invariants of the quadratic form .
The quadratic form is positive definite if (similarly, negative definite if ) for every nonzero vector . When assumes both positive and negative values, is an isotropic quadratic form. The theorems of Jacobi and Sylvester show that any positive definite quadratic form in variables can be brought to the sum of squares by a suitable invertible linear transformation: geometrically, there is only one positive definite real quadratic form of every dimension. Its isometry group is a compact orthogonal group . This stands in contrast with the case of isotropic forms, when the corresponding group, the indefinite orthogonal group , is non-compact. Further, the isometry groups of and are the same (, but the associated Clifford algebras (and hence pin groups) are different.
Definitions
A quadratic form over a field is a map from a finite-dimensional -vector space to such that for all , and the function is bilinear.
More concretely, an -ary quadratic form over a field is a homogeneous polynomial of degree 2 in variables with coefficients in :
This formula may be rewritten using matrices: let be the column vector with components , ..., and be the matrix over whose entries are the coefficients of . Then
A vector is a null vector if .
Two -ary quadratic forms and over are equivalent if there exists a nonsingular linear transformation such that
Let the characteristic of be different from 2. The coefficient matrix of may be replaced by the symmetric matrix with the same quadratic form, so it may be assumed from the outset that is symmetric. Moreover, a symmetric matrix is uniquely determined by the corresponding quadratic form. Under an equivalence , the symmetric matrix of and the symmetric matrix of are related as follows:
The associated bilinear form of a quadratic form is defined by
Thus, is a symmetric bilinear form over with matrix . Conversely, any symmetric bilinear form defines a quadratic form
and these two processes are the inverses of each other. As a consequence, over a field of characteristic not equal to 2, the theories of symmetric bilinear forms and of quadratic forms in variables are essentially the same.
Quadratic space
Given an -dimensional vector space over a field , a quadratic form on is a function that has the following property: for some basis, the function that maps the coordinates of to is a quadratic form. In particular, if with its standard basis, one has
The change of basis formulas show that the property of being a quadratic form does not depend on the choice of a specific basis in , although the quadratic form depends on the choice of the basis.
A finite-dimensional vector space with a quadratic form is called a quadratic space.
The map is a homogeneous function of degree 2, which means that it has the property that, for all in and in :
When the characteristic of is not 2, the bilinear map over is defined:
This bilinear form is symmetric. That is, for all , in , and it determines : for all in .
When the characteristic of is 2, so that 2 is not a unit, it is still possible to use a quadratic form to define a symmetric bilinear form . However, can no longer be recovered from this in the same way, since for all (and is thus alternating). Alternatively, there always exists a bilinear form (not in general either unique or symmetric) such that .
The pair consisting of a finite-dimensional vector space over and a quadratic map from to is called a quadratic space, and as defined here is the associated symmetric bilinear form of . The notion of a quadratic space is a coordinate-free version of the notion of quadratic form. Sometimes, is also called a quadratic form.
Two -dimensional quadratic spaces and are isometric if there exists an invertible linear transformation (isometry) such that
The isometry classes of -dimensional quadratic spaces over correspond to the equivalence classes of -ary quadratic forms over .
Generalization
Let be a commutative ring, be an -module, and be an -bilinear form. A mapping is the associated quadratic form of , and is the polar form of .
A quadratic form may be characterized in the following equivalent ways:
There exists an -bilinear form such that is the associated quadratic form.
for all and , and the polar form of is -bilinear.
Related concepts
Two elements and of are called orthogonal if . The kernel of a bilinear form consists of the elements that are orthogonal to every element of . is non-singular if the kernel of its associated bilinear form is . If there exists a non-zero in such that , the quadratic form is isotropic, otherwise it is definite. This terminology also applies to vectors and subspaces of a quadratic space. If the restriction of to a subspace of is identically zero, then is totally singular.
The orthogonal group of a non-singular quadratic form is the group of the linear automorphisms of that preserve : that is, the group of isometries of into itself.
If a quadratic space has a product so that is an algebra over a field, and satisfies
then it is a composition algebra.
Equivalence of forms
Every quadratic form in variables over a field of characteristic not equal to 2 is equivalent to a diagonal form
Such a diagonal form is often denoted by .
Classification of all quadratic forms up to equivalence can thus be reduced to the case of diagonal forms.
Geometric meaning
Using Cartesian coordinates in three dimensions, let , and let be a symmetric 3-by-3 matrix. Then the geometric nature of the solution set of the equation depends on the eigenvalues of the matrix .
If all eigenvalues of are non-zero, then the solution set is an ellipsoid or a hyperboloid. If all the eigenvalues are positive, then it is an ellipsoid; if all the eigenvalues are negative, then it is an imaginary ellipsoid (we get the equation of an ellipsoid but with imaginary radii); if some eigenvalues are positive and some are negative, then it is a hyperboloid.
If there exist one or more eigenvalues , then the shape depends on the corresponding . If the corresponding , then the solution set is a paraboloid (either elliptic or hyperbolic); if the corresponding , then the dimension degenerates and does not come into play, and the geometric meaning will be determined by other eigenvalues and other components of . When the solution set is a paraboloid, whether it is elliptic or hyperbolic is determined by whether all other non-zero eigenvalues are of the same sign: if they are, then it is elliptic; otherwise, it is hyperbolic.
Integral quadratic forms
Quadratic forms over the ring of integers are called integral quadratic forms, whereas the corresponding modules are quadratic lattices (sometimes, simply lattices). They play an important role in number theory and topology.
An integral quadratic form has integer coefficients, such as ; equivalently, given a lattice in a vector space (over a field with characteristic 0, such as or ), a quadratic form is integral with respect to if and only if it is integer-valued on , meaning if .
This is the current use of the term; in the past it was sometimes used differently, as detailed below.
Historical use
Historically there was some confusion and controversy over whether the notion of integral quadratic form should mean:
twos in the quadratic form associated to a symmetric matrix with integer coefficients
twos out a polynomial with integer coefficients (so the associated symmetric matrix may have half-integer coefficients off the diagonal)
This debate was due to the confusion of quadratic forms (represented by polynomials) and symmetric bilinear forms (represented by matrices), and "twos out" is now the accepted convention; "twos in" is instead the theory of integral symmetric bilinear forms (integral symmetric matrices).
In "twos in", binary quadratic forms are of the form , represented by the symmetric matrix
This is the convention Gauss uses in Disquisitiones Arithmeticae.
In "twos out", binary quadratic forms are of the form , represented by the symmetric matrix
Several points of view mean that twos out has been adopted as the standard convention. Those include:
better understanding of the 2-adic theory of quadratic forms, the 'local' source of the difficulty;
the lattice point of view, which was generally adopted by the experts in the arithmetic of quadratic forms during the 1950s;
the actual needs for integral quadratic form theory in topology for intersection theory;
the Lie group and algebraic group aspects.
Universal quadratic forms
An integral quadratic form whose image consists of all the positive integers is sometimes called universal. Lagrange's four-square theorem shows that is universal. Ramanujan generalized this and found 54 multisets that can each generate all positive integers, namely,
There are also forms whose image consists of all but one of the positive integers. For example, has 15 as the exception. Recently, the 15 and 290 theorems have completely characterized universal integral quadratic forms: if all coefficients are integers, then it represents all positive integers if and only if it represents all integers up through 290; if it has an integral matrix, it represents all positive integers if and only if it represents all integers up through 15.
| Mathematics | Abstract algebra | null |
251489 | https://en.wikipedia.org/wiki/Ironclad%20warship | Ironclad warship | An ironclad was a steam-propelled warship protected by steel or iron armor constructed from 1859 to the early 1890s. The ironclad was developed as a result of the vulnerability of wooden warships to explosive or incendiary shells. The first ironclad battleship, , was launched by the French Navy in November 1859, narrowly preempting the British Royal Navy. However, Britain built the first completely iron-hulled warships.
They were first used in warfare in 1862 during the American Civil War, when ironclads operated against wooden ships and, in a historic confrontation, against each other at the Battle of Hampton Roads in Virginia. Their performance demonstrated that the ironclad had replaced the unarmored ship of the line as the most powerful warship afloat. Ironclad gunboats became very successful in the American Civil War.
Ironclads were designed for several uses, including as high-seas battleships, long-range cruisers, and coastal defense ships. Rapid development of warship design in the late 19th century transformed the ironclad from a wooden-hulled vessel that carried sails to supplement its steam engines into the steel-built, turreted battleships, and cruisers familiar in the 20th century. This change was pushed forward by the development of heavier naval guns, more sophisticated steam engines, and advances in ferrous metallurgy that made steel shipbuilding possible.
The quick pace of change meant that many ships were obsolete almost as soon as they were finished and that naval tactics were in a state of flux. Many ironclads were built to make use of the naval ram, the torpedo, or sometimes both (as in the case with smaller ships and later torpedo boats), which several naval designers considered the important weapons of naval combat. There is no clear end to the ironclad period, but toward the end of the 1890s, the term ironclad dropped out of use. New ships were increasingly constructed to a standard pattern and designated as battleships or armored cruisers.
Development
The ironclad became technically feasible and tactically necessary because of developments in shipbuilding in the first half of the 19th century. According to naval historian J. Richard Hill: "The (ironclad) had three chief characteristics: a metal-skinned hull, steam propulsion and a main armament of guns capable of firing explosive shells. It is only when all three characteristics are present that a fighting ship can properly be called an ironclad." Each of these developments was introduced separately in the decade before the first ironclads.
Steam propulsion
In the 18th and early 19th centuries, fleets had relied on two types of major warship, the ship of the line and the frigate. The first major change to these types was the introduction of steam power for propulsion. While paddle steamer warships had been used from the 1830s onward, steam propulsion only became suitable for major warships after the adoption of the screw propeller in the 1840s.
Steam-powered screw frigates were built in the mid-1840s, and at the end of the decade the French Navy introduced steam power to its line of battle. Napoleon III's ambition to gain greater influence in Europe required a sustained challenge to the British at sea. The first purpose-built steam battleship was the 90-gun in 1850. Napoléon was armed as a conventional ship-of-the-line, but her steam engines could give her a speed of , regardless of the wind conditions: a potentially decisive advantage in a naval engagement.
The introduction of the steam ship-of-the-line led to a building competition between France and Britain. Eight sister ships to Napoléon were built in France over a period of ten years, but the United Kingdom soon managed to take the lead in production. Altogether, France built ten new wooden steam battleships and converted 28 from older ships of the line, while the United Kingdom built 18 and converted 41.
Explosive shells
The era of the wooden steam ship-of-the-line was brief, because of new, more powerful naval guns. In the 1820s and 1830s, warships began to mount increasingly heavy guns, replacing 18- and 24-pounder guns with 32-pounders on sailing ships-of-the-line and introducing 68-pounders on steamers. Then, the first shell guns firing explosive shells were introduced following their development by the French Général Henri-Joseph Paixhans. By the 1840s they were part of the standard armament for naval powers including the French Navy, Royal Navy, Imperial Russian Navy and United States Navy.
It is often held that the power of explosive shells to smash wooden hulls, as demonstrated by the Russian destruction of an Ottoman squadron at the Battle of Sinop, spelled the end of the wooden-hulled warship. The more practical threat to wooden ships was from conventional cannon firing red-hot shot, which could lodge in the hull and cause a fire or ammunition explosion. Some navies even experimented with hollow shot filled with molten metal for extra incendiary power.
Iron armor
The use of wrought iron instead of wood as the primary material of ships' hulls began in the 1830s; the first "warship" with an iron hull was the gunboat Nemesis, built by Jonathan Laird of Birkenhead for the East India Company in 1839. There followed, also from Laird, the first full-sized warship with a metal hull, the 1842 steam frigate Guadalupe for the Mexican Navy. The latter ship performed well during the Naval Battle of Campeche, with her captain reporting that he thought that there were fewer iron splinters from Guadalupes hull than from a wooden hull.
Encouraged by the positive reports of the iron hulls of those ships in combat, the Admiralty ordered a series of experiments to evaluate what happened when thin iron hulls were struck by projectiles, both solid shot and hollow shells, beginning in 1845 and lasting through 1851. Critics like Lieutenant-general Sir Howard Douglas believed that the splinters from the hull were even more dangerous than those from wooden hulls and the tests partially confirmed this belief. What was ignored was that of wood backing the iron would stop most of the splinters from penetrating and that relatively thin plates of iron backed by the same thickness of wood would generally cause shells to split open and fail to detonate. One factor in the performance of wrought iron during these tests that was not understood by metallurgists of the day was that wrought iron begins to become brittle at temperatures below . Many of the tests were conducted at temperatures below this while the battles were fought in tropical climates.
The early experimental results seemed to support the critics and party politics came into play as the Whig First Russell ministry replaced the Tory Second Peel Ministry in 1846. The new administration sided with the critics and ordered that the four iron-hulled propeller frigates ordered by the Tories be converted into troopships. No iron warships would be ordered until the beginning of the Crimean War in 1854.
Following the demonstration of the power of explosive shells against wooden ships at the Battle of Sinop, and fearing that his own ships would be vulnerable to the Paixhans guns of Russian fortifications in the Crimean War, Emperor Napoleon III ordered the development of light-draft floating batteries, equipped with heavy guns and protected by heavy armor. Experiments made during the first half of 1854 proved highly satisfactory, and on 17 July 1854, the French communicated to the British Government that a solution had been found to make gun-proof vessels and that plans would be communicated. After tests in September 1854, the British Admiralty agreed to build five armored floating batteries on the French plans.
The French floating batteries were deployed in 1855 as a supplement to the wooden steam battle fleet in the Crimean War. The role of the battery was to assist unarmored mortar and gunboats bombarding shore fortifications. The French used three of their ironclad batteries (Lave, Tonnante and Dévastation) in 1855 against the defenses at the Battle of Kinburn on the Black Sea, where they were effective against Russian shore defences. They would later be used again during the Italian war in the Adriatic in 1859. The British floating batteries and arrived too late to participate to the action at Kinburn. The British planned to use theirs in the Baltic Sea against the well-fortified Russian naval base at Kronstadt.
The batteries have a claim to the title of the first ironclad warships but they were capable of only under their own power: they operated under their own power at the Battle of Kinburn, but had to be towed for long-range transit. They were also arguably marginal to the work of the navy. The brief success of the floating ironclad batteries convinced France to begin work on armored warships for their battlefleet.
Early ironclad ships and battles
By the end of the 1850s it was clear that France was unable to match British building of steam warships, and to regain the strategic initiative a dramatic change was required. The result was the first ocean-going ironclad, , begun in 1857 and launched in 1859. Gloires wooden hull was modelled on that of a steam ship of the line, reduced to one deck, and sheathed in iron plates thick. She was propelled by a steam engine, driving a single screw propeller for a speed of . She was armed with thirty-six rifled guns. France proceeded to construct 16 ironclad warships, including two sister ships to Gloire, and the only two-decked broadside ironclads ever built, and .
The Royal Navy had not been keen to sacrifice its advantage in steam ships of the line, but was determined that the first British ironclad would outmatch the French ships in every respect, particularly speed. A fast ship would have the advantage of being able to choose a range of engagement that could make her invulnerable to enemy fire. The British specification was more a large, powerful frigate than a ship-of-the-line. The requirement for speed meant a very long vessel, which had to be built from iron. The result was the construction of two s; and . The ships had a successful design, though there were necessarily compromises between 'sea-keeping', strategic range and armor protection. Their weapons were more effective than those of Gloire, and with the largest set of steam engines yet fitted to a ship, they could steam at 14.3 knots (26.5 km/h). Yet the Gloire and her sisters had full iron-armor protection along the waterline and the battery itself.
The British Warrior and Black Prince (but also the smaller Defence and Resistance) were obliged to concentrate their armor in a central "citadel" or "armoured box", leaving many main deck guns and the fore and aft sections of the vessel unprotected. The use of iron in the construction of Warrior also came with some drawbacks; iron hulls required more regular and intensive repairs than wooden hulls, and iron was more susceptible to fouling by marine life.
By 1862, navies across Europe had adopted ironclads. Britain and France each had sixteen either completed or under construction, though the British vessels were larger. Austria, Italy, Russia, and Spain were also building ironclads. However, the first battles using the new ironclad ships took place during the American Civil War, between Union and Confederate ships in 1862. These were markedly different from the broadside-firing, masted designs of Gloire and Warrior. The clash of the Italian and Austrian fleets at the Battle of Lissa (1866), also had an important influence on the development of ironclad design.
First battles between ironclads: the U.S. Civil War
The first use of ironclads in combat came in the U.S. Civil War. The U.S. Navy at the time the war broke out had no ironclads, its most powerful ships being six unarmored steam-powered frigates. Since the bulk of the Navy remained loyal to the Union, the Confederacy sought to gain advantage in the naval conflict by acquiring modern armored ships. In May 1861, the Confederate Congress appropriated $2 million dollars for the purchase of ironclads from overseas, and in July and August 1861 the Confederacy started work on construction and converting wooden ships.
On 12 October 1861, became the first ironclad to enter combat, when she fought Union warships on the Mississippi during the Battle of the Head of Passes. She had been converted from a commercial vessel in New Orleans for river and coastal fighting. In February 1862, the larger joined the Confederate Navy, having been rebuilt at Norfolk. Constructed on the hull of , Virginia originally was a conventional warship made of wood, but she was converted into an iron-covered casemate ironclad gunship, when she entered the Confederate Navy. By this time, the Union had completed seven ironclad gunboats of the , and was about to complete , an innovative design proposed by the Swedish inventor John Ericsson. The Union was also building a large armored frigate, , and the smaller .
The first battle between ironclads happened on 9 March 1862, as the armored Monitor was deployed to protect the Union's wooden fleet from the ironclad ram Virginia and other Confederate warships. In this engagement, the second day of the Battle of Hampton Roads, the two ironclads tried to ram one another while shells bounced off their armor. The battle attracted attention worldwide, making it clear that the wooden warship was now out of date, with the ironclads destroying them easily.
The Civil War saw more ironclads built by both sides, and they played an increasing role in the naval war alongside the unarmored warships, commerce raiders and blockade runners. The Union built a large fleet of fifty monitors modeled on their namesake. The Confederacy built ships designed as smaller versions of Virginia, many of which saw action, but their attempts to buy ironclads overseas were frustrated as European nations confiscated ships being built for the Confederacy – especially in Russia, the only country to openly support the Union through the war. Only CSS Stonewall was completed, and she arrived in Cuban waters just in time for the end of the war.
Through the remainder of the war, ironclads saw action in the Union's attacks on Confederate ports. Seven Union monitors, including , as well as two other ironclads, the armored frigate New Ironsides and a light-draft , participated in the failed attack on Charleston; one was sunk. Two small ironclads, and participated in the defense of the harbor. For the later attack at Mobile Bay, the Union assembled four monitors as well as 11 wooden ships, facing the , the Confederacy's most powerful ironclad, and three gunboats.
On the western front, the Union built a formidable force of river ironclads, beginning with several converted riverboats and then contracting engineer James Eads of St. Louis, Missouri to build the City-class ironclads. These excellent ships were built with twin engines and a central paddle wheel, all protected by an armored casemate. They had a shallow draft, allowing them to journey up smaller tributaries, and were very well suited for river operations. Eads also produced monitors for use on the rivers, the first two of which differed from the ocean-going monitors in that they contained a paddle wheel ( and ).
The Union ironclads played an important role in the Mississippi and tributaries by providing tremendous fire upon Confederate forts, installations and vessels with relative impunity to enemy fire. They were not as heavily armored as the ocean-going monitors of the Union, but they were adequate for their intended use. More Western Flotilla Union ironclads were sunk by torpedoes (mines) than by enemy fire, and the most damaging fire for the Union ironclads was from shore installations, not Confederate vessels.
Lissa: first fleet battle
The first fleet battle, and the first ocean battle, involving ironclad warships was the Battle of Lissa in 1866. Waged between the Austrian and Italian navies, the battle pitted combined fleets of wooden frigates and corvettes and ironclad warships on both sides in the largest naval battle between the battles of Navarino and Tsushima.
The Italian fleet consisted of 12 ironclads and a similar number of wooden warships, escorting transports which carried troops intending to land on the Adriatic island of Lissa. Among the Italian ironclads were seven broadside ironclad frigates, four smaller ironclads, and the newly built – a double-turreted ram. Opposing them, the Austrian navy had seven ironclad frigates.
The Austrians believed their ships to have less effective guns than their enemy, so decided to engage the Italians at close range and ram them. The Austrian fleet formed into an arrowhead formation with the ironclads in the first line, charging at the Italian ironclad squadron. In the melée which followed both sides were frustrated by the lack of damage inflicted by guns, and by the difficulty of ramming—nonetheless, the effective ramming attack being made by the Austrian flagship against the Italian attracted great attention in following years.
The superior Italian fleet lost its two ironclads, and , while the Austrian unarmored screw two-decker remarkably survived close actions with four Italian ironclads. The battle ensured the popularity of the ram as a weapon in European ironclads for many years, and the victory won by Austria established it as the predominant naval power in the Adriatic.
The battles of the American Civil War and at Lissa were very influential on the designs and tactics of the ironclad fleets that followed. In particular, it taught a generation of naval officers the (ultimately erroneous) lesson that ramming was the best way to sink enemy ironclads.
Armament and tactics
The adoption of iron armor meant that the traditional naval armament of dozens of light cannon became useless, since their shot would bounce off an armored hull. To penetrate armor, increasingly heavy guns were mounted on ships; nevertheless, the view that ramming was the only way to sink an ironclad became widespread. The increasing size and weight of guns also meant a movement away from the ships mounting many guns broadside, in the manner of a ship-of-the-line, towards a handful of guns in turrets for all-round fire.
Ram craze
From the 1860s to the 1880s many naval designers believed that the ram was again a vital weapon in naval warfare. With steam power freeing ships from the wind, iron construction increasing their structural strength, and armor making them invulnerable to shellfire, the ram seemed to offer the opportunity to strike a decisive blow.
The scant damage inflicted by the guns of Monitor and Virginia at Hampton Roads and the spectacular but lucky success of the Austrian flagship SMS Erzherzog Ferdinand Max sinking the Italian Re d'Italia at Lissa gave strength to the ramming craze. From the early 1870s to early 1880s most British naval officers thought that guns were about to be replaced as the main naval armament by the ram. Those who noted the tiny number of ships that had actually been sunk by ramming struggled to be heard.
The revival of ramming had a significant effect on naval tactics. Since the 17th century the predominant tactic of naval warfare had been the line of battle, where a fleet formed a long line to give it the best fire from its broadside guns. This tactic was totally unsuited to ramming, and the ram threw fleet tactics into disarray. The question of how an ironclad fleet should deploy in battle to make best use of the ram was never tested in battle, and if it had been, combat might have shown that rams could only be used against ships which were already stopped dead in the water.
The ram finally fell out of favor in the 1880s, as the same effect could be achieved with a torpedo, with less vulnerability to quick-firing guns.
Development of naval guns
The armament of ironclads tended to become concentrated in a small number of powerful guns capable of penetrating the armor of enemy ships at range; calibre and weight of guns increased markedly to achieve greater penetration. Throughout the ironclad era navies also grappled with the complexities of rifled versus smoothbore guns and breech-loading versus muzzle-loading.
carried a mixture of 110-pounder breech-loading rifles and more traditional 68-pounder smoothbore guns. Warrior highlighted the challenges of picking the right armament; the breech-loaders she carried, designed by Sir William Armstrong, were intended to be the next generation of heavy armament for the Royal Navy, but were shortly withdrawn from service.
Breech-loading guns seemed to offer important advantages. A breech-loader could be reloaded without moving the gun, a lengthy process particularly if the gun then needed to be re-aimed. Warriors Armstrong guns also had the virtue of being lighter than an equivalent smoothbore and, because of their rifling, more accurate. Nonetheless, the design was rejected because of problems which plagued breech-loaders for decades.
The weakness of the breech-loader was the obvious problem of sealing the breech. All guns are powered by the explosive conversion of a solid propellant into gas. This explosion propels the shot or shell out of the front of the gun, but also imposes great stresses on the gun-barrel. If the breech—which experiences some of the greatest forces in the gun—is not entirely secure, then there is a risk that either gas will discharge through the breech or that the breech will break. This in turn reduces the muzzle velocity of the weapon and can also endanger the gun crew. Warriors Armstrong guns suffered from both problems; the shells were unable to penetrate the armor of Gloire, while sometimes the screw which closed the breech flew backwards out of the gun on firing. Similar problems were experienced with the breech-loading guns which became standard in the French and German navies.
These problems influenced the British to equip ships with muzzle-loading weapons of increasing power until the 1880s. After a brief introduction of the 100-pounder or smoothbore Somerset Gun, which weighed , the Admiralty introduced 7-inch (178 mm) rifled guns, weighing . These were followed by a series of increasingly mammoth weapons—guns weighing , , , and finally , with caliber increasing from 8 inches (203 mm) to 16 inches (406 mm).
The decision to retain muzzle-loaders until the 1880s has been criticized by historians. However, at least until the late 1870s, the British muzzle-loaders had superior performance in terms of both range and rate of fire than the French and Prussian breech-loaders, which suffered from the same problems as the first Armstrong guns.
From 1875 onwards, the balance between breech- and muzzle-loading changed. Captain de Bange invented a method of reliably sealing a breech, adopted by the French in 1873. Just as compellingly, the growing size of naval guns and consequently, their ammunition, made muzzle-loading much more complicated. With guns of such size there was no prospect of hauling in the gun for reloading, or even reloading by hand, and complicated hydraulic systems were required for reloading the gun outside the turret without exposing the crew to enemy fire. In 1882, the 81-ton, 16-inch guns of fired only once every 11 minutes while bombarding Alexandria during the Urabi Revolt. The , 450 mm (17.72 inch) guns of the could each fire a round every 15 minutes.
In the Royal Navy, the switch to breech-loaders was finally made in 1879; as well as the significant advantages in terms of performance, opinion was swayed by an explosion on board caused by a gun being double-loaded, a problem which could only happen with a muzzle-loading gun.
The caliber and weight of guns could only increase so far. The larger the gun, the slower it would be to load, the greater the stresses on the ship's hull, and the less the stability of the ship. The size of the gun peaked in the 1880s, with some of the heaviest calibers of gun ever used at sea. carried two 16.25-inch (413 mm) breech-loading guns, each weighing . A few years afterwards, the Italians used 450 mm (17.72 inch) muzzle-loading guns on the Duilio class ships. One consideration which became more acute was that even from the original Armstrong models, following the Crimean War, range and hitting power far exceeded simple accuracy, especially at sea where the slightest roll or pitch of the vessel as 'floating weapons-platform' could negate the advantage of rifling. American ordnance experts accordingly preferred smoothbore monsters whose round shot could at least 'skip' along the surface of the water. Actual effective combat ranges, they had learned during the Civil War, were comparable to those in the Age of Sail—though a vessel could now be smashed to pieces in only a few rounds. Smoke and the general chaos of battle only added to the problem. As a result, many naval engagements in the 'Age of the Ironclad' were still fought at ranges within easy eyesight of their targets, and well below the maximum reach of their ships' guns.
Another method of increasing firepower was to vary the projectile fired or the nature of the propellant. Early ironclads used black powder, which expanded rapidly after combustion; this meant cannons had relatively short barrels, to prevent the barrel itself slowing the shell. The sharpness of the black powder explosion also meant that guns were subjected to extreme stress. One important step was to press the powder into pellets, allowing a slower, more controlled explosion and a longer barrel. A further step forward was the introduction of chemically different brown powder which combusted more slowly again. It also put less stress on the insides of the barrel, allowing guns to last longer and to be manufactured to tighter tolerances.
The development of smokeless powder, based on nitroglycerine or nitrocellulose, by the French inventor Paul Vielle in 1884 was a further step allowing smaller charges of propellant with longer barrels. The guns of the pre-Dreadnought battleships of the 1890s tended to be smaller in caliber compared to the ships of the 1880s, most often 12 in (305 mm), but progressively grew in length of barrel, making use of improved propellants to gain greater muzzle velocity.
The nature of the projectiles also changed during the ironclad period. Initially, the best armor-piercing projectile was a solid cast-iron shot. Later, shot of chilled iron, a harder iron alloy, gave better armor-piercing qualities. Eventually the armor-piercing shell was developed.
Positioning of armament
Broadside ironclads
The first British, French and Russian ironclads, in a logical development of warship design from the long preceding era of wooden ships of the line, carried their weapons in a single line along their sides and so were called "broadside ironclads". Both and were examples of this type. Because their armor was so heavy, they could only carry a single row of guns along the main deck on each side rather than a row on each deck.
A significant number of broadside ironclads were built in the 1860s, principally in Britain and France, but in smaller numbers by other powers including Italy, Austria, Russia and the United States. The advantages of mounting guns on both broadsides was that the ship could engage more than one adversary at a time, and the rigging did not impede the field of fire.
Broadside armament also had disadvantages, which became more serious as ironclad technology developed. Heavier guns to penetrate ever-thicker armor meant that fewer guns could be carried. Furthermore, the adoption of ramming as an important tactic meant the need for ahead and all-round fire. These problems led to broadside designs being superseded by designs that gave greater all-round fire, which included central-battery, turret, and barbette designs.
Turrets, batteries, and barbettes
There were two main design alternatives to the broadside. In one design, the guns were placed in an armored casemate amidships: this arrangement was called the 'box-battery' or 'center-battery'. In the other, the guns could be placed on a rotating platform to give them a broad field of fire; when fully armored, this arrangement was called a turret and when partially armored or unarmored, a barbette.
The centre-battery was the simpler and, during the 1860s and 1870s, the more popular method. Concentrating guns amidships meant the ship could be shorter and handier than a broadside type. The first full-scale center-battery ship was of 1865; the French laid down centre-battery ironclads in 1865 which were not completed until 1870. Centre-battery ships often, but not always, had a recessed freeboard enabling some of their guns to fire directly ahead.
The turret was first used in naval combat on in 1862, with a type of turret designed by the Swedish engineer John Ericsson. A competing turret design was proposed by the British inventor Cowper Coles with a prototype of this installed on in 1861 for testing and evaluation purposes. Ericsson's turret turned on a central spindle, and Coles's turned on a ring of bearings. Turrets offered the maximum arc of fire for the guns, but there were significant problems with their use in the 1860s. The fire arc of a turret would be considerably limited by masts and rigging, so they were unsuited to use on the earlier ocean-going ironclads. The second problem was that turrets were extremely heavy. Ericsson was able to offer the heaviest possible turret (guns and armor protection) by deliberately designing a ship with very low freeboard. The weight thus saved from having a high broadside above the waterline was diverted to actual guns and armor. Low freeboard, however, also meant a smaller hull and therefore a smaller capacity for coal storage—and therefore range of the vessel. In many respects, the turreted, low-freeboard Monitor and the broadside sailor HMS Warrior represented two opposite extremes in what an 'Ironclad' was all about. The most dramatic attempt to compromise these two extremes, or 'squaring this circle', was designed by Captain Cowper Phipps Coles: . It was a dangerously low freeboard turret ship, which nevertheless carried a full rig of sail and subsequently capsized not long after her launch in 1870. Her half-sister was restricted to firing from her turrets only on the port and starboard beams. The third Royal Navy ship to combine turrets and masts was of 1876, which carried two turrets on either side of the center-line, allowing both to fire fore, aft and broadside.
A lighter alternative to the turret, particularly popular with the French navy, was the barbette. These were fixed armored towers which held a gun on a turntable. The crew was sheltered from direct fire, but vulnerable to plunging fire, for instance from shore emplacements. The barbette was lighter than the turret, needing less machinery and no roof armor. Some barbettes were stripped of their armor plate to reduce the top-weight of their ships. The barbette became widely adopted in the 1880s, and with the addition of an armored 'gun-house', transformed into the turrets of the pre-dreadnought battleships.
Torpedoes
The ironclad age saw the development of explosive torpedoes as naval weapons, which helped complicate the design and tactics of ironclad fleets. The first torpedoes were static mines, used extensively in the American Civil War. That conflict also saw the development of the spar torpedo, an explosive charge pushed against the hull of a warship by a small boat. For the first time, a large warship faced a serious threat from a smaller one—and given the relative inefficiency of shellfire against ironclads, the threat from the spar torpedo was taken seriously. The U.S. Navy converted four of its monitors to become turretless armored spar-torpedo vessels while under construction in 1864–1865, but these vessels never saw action. Another proposal, the towed or 'Harvey' torpedo, involved an explosive on a line or outrigger; either to deter a ship from ramming or to make a torpedo attack by a boat less suicidal.
A more practical and influential weapon was the self-propelled or Whitehead torpedo. Invented in 1868 and deployed in the 1870s, it formed part of the armament of ironclads of the 1880s like HMS Inflexible and the Italian . The ironclad's vulnerability to the torpedo was a key part of the critique of armored warships made by the school of naval thought; it appeared that any ship armored enough to prevent destruction by gunfire would be slow enough to be easily caught by torpedo. In practice, however, the Jeune Ecole was only briefly influential and the torpedo formed part of the confusing mixture of weapons possessed by ironclads.
Armor and construction
The first ironclads were built on wooden or iron hulls, and protected by wrought iron armor backed by thick wooden planking. Ironclads were still being built with wooden hulls into the 1870s.
Hulls: iron, wood, and steel
Using wrought iron construction for warships offered advantages for the engineering of the hull. However, unarmored iron had many military disadvantages, and offered technical problems which kept wooden hulls in use for many years, particularly for long-range cruising warships.
Iron ships had first been proposed for military use in the 1820s. In the 1830s and 1840s, France, Britain and the United States had all experimented with iron-hulled but unarmored gunboats and frigates. However, the iron-hulled frigate was abandoned by the end of the 1840s, because iron hulls were more vulnerable to solid shot; iron was more brittle than wood, and iron frames more likely to fall out of shape than wood.
The unsuitability of unarmored iron for warship hulls meant that iron was only adopted as a building material for battleships when protected by armor. However, iron gave the naval architect many advantages. Iron allowed larger ships and more flexible design, for instance the use of watertight bulkheads on the lower decks. Warrior, built of iron, was longer and faster than the wooden-hulled Gloire. Iron could be produced to order and used immediately, in contrast to the need to give wood a long period of seasoning. And, given the large quantities of wood required to build a steam warship and the falling cost of iron, iron hulls were increasingly cost-effective. The main reason for the French use of wooden hulls for the ironclad fleet built in the 1860s was that the French iron industry could not supply enough, and the main reason why Britain built its handful of wooden-hulled ironclads was to make best use of hulls already started and wood already bought.
Wooden hulls continued to be used for long-range and smaller ironclads, because iron nevertheless had a significant disadvantage. Iron hulls suffered quick fouling by marine life, slowing the ships down—manageable for a European battlefleet close to dry docks, but a difficulty for long-range ships. The only solution was to sheath the iron hull first in wood and then in copper, a laborious and expensive process which made wooden construction remain attractive. Iron and wood were to some extent interchangeable: the Japanese and ordered in 1875 were sister-ships, but one was built of iron and the other of composite construction.
After 1872, steel started to be introduced as a material for construction. Compared to iron, steel allows for greater structural strength for a lower weight. The French Navy led the way with the use of steel in its fleet, starting with the , laid down in 1873 and launched in 1876. Redoutable nonetheless had wrought iron armor plate, and part of her exterior hull was iron rather than steel.
Even though Britain led the world in steel production, the Royal Navy was slow to adopt steel warships. The Bessemer process for steel manufacture produced too many imperfections for large-scale use on ships. French manufacturers used the Siemens-Martin process to produce adequate steel, but British technology lagged behind. The first all-steel warships built by the Royal Navy were the dispatch vessels Iris and Mercury, laid down in 1875 and 1876.
Armor and protection schemes
Iron-built ships used wood as part of their protection scheme. HMS Warrior was protected by 4.5 in (114 mm) of wrought iron backed by 15 in (381 mm) of teak, the strongest shipbuilding wood. The wood played two roles, preventing spalling and also preventing the shock of a hit damaging the structure of the ship. Later, wood and iron were combined in 'sandwich' armor, for instance in HMS Inflexible.
Steel was also an obvious material for armor. It was tested in the 1860s, but the steel of the time was too brittle and disintegrated when struck by shells. Steel became practical to use when a way was found to fuse steel onto wrought iron plates, giving a form of compound armor. This compound armor was used by the British in ships built from the late 1870s, first for turret armor (starting with HMS Inflexible) and then for all armor (starting with of 1882). The French and German navies adopted the innovation almost immediately, with licenses being given for the use of the 'Wilson System' of producing fused armor.
The first ironclads to have all-steel armor were the two ships of the Duilio class. Though the ships were laid down in 1873 their armor was not purchased from France until 1877. The French navy decided in 1880 to adopt compound armor for its fleet, but found it limited in supply, so from 1884 the French navy was using steel armor. Britain stuck to compound armor until 1889.
The ultimate ironclad armor was case hardened nickel-steel. In 1890, the U.S. Navy tested steel armor hardened by the Harvey process and found it superior to compound armor. For several years 'Harvey steel' was the state of the art, produced in the U.S., France, Germany, Britain, Austria and Italy. In 1894, the German firm Krupp developed gas cementing, which further hardened steel armor. The German , laid down in 1895, was the first ship to benefit from the new 'Krupp armor' and the new armor was quickly adopted; the Royal Navy using it from , laid down in 1896. By 1901 almost all new battleships used Krupp armor, though the U.S. continued to use Harvey armor alongside until the end of the decade.
The equivalent strengths of the different armor plates was as follows: 15 in (381 mm) of wrought iron was equivalent to 12 in (305 mm) of either plain steel or compound iron and steel armor, and to 7.75 in (197 mm) of Harvey armor or 5.75 in (146 mm) of Krupp armor.
Ironclad construction also prefigured the later debate in battleship design between tapering and 'all-or-nothing' armor design. Warrior was only semi-armored, and could have been disabled by hits on the bow and stern. As the thickness of armor grew to protect ships from the increasingly heavy guns, the area of the ship which could be fully protected diminished. Inflexibles armor protection was largely limited to the central citadel amidships, protecting boilers and engines, turrets and magazines, and little else. An ingenious arrangement of cork-filled compartments and watertight bulkheads was intended to keep her stable and afloat in the event of damage to her un-armored sections.
Propulsion: steam and sail
The first ocean-going ironclads carried masts and sails like their wooden predecessors, and these features were only gradually abandoned. Early steam engines were inefficient; the wooden steam fleet of the Royal Navy could only carry "5 to 9 days coal", and the situation was similar with the early ironclads. Warrior also illustrates two design features which aided hybrid propulsion; she had retractable screws to reduce drag while under sail (though in practice the steam engine was run at a low throttle), and a telescopic funnel which could be folded down to the deck level.
Ships designed for coastal warfare, like the floating batteries of the Crimea, or and her sisters, dispensed with masts from the beginning. The British , started in 1869, was the first large, ocean-going ironclad to dispense with masts. Her principal role was for combat in the English Channel and other European waters; while her coal supplies gave her enough range to cross the Atlantic, she would have had little endurance on the other side of the ocean. The Devastation and the similar ships commissioned by the British and Russian navies in the 1870s were the exception rather than the rule. Most ironclads of the 1870s retained masts, and only the Italian navy, which during that decade was focused on short-range operations in the Adriatic, built consistently mastless ironclads.
During the 1860s, steam engines improved with the adoption of double-expansion steam engines, which used 30–40% less coal than earlier models. The Royal Navy decided to switch to the double-expansion engine in 1871, and by 1875 they were widespread. However, this development alone was not enough to herald the end of the mast. Whether this was due to a conservative desire to retain sails, or was a rational response to the operational and strategic situation, is a matter of debate. A steam-only fleet would require a network of coaling stations worldwide, which would need to be fortified at great expense to stop them falling into enemy hands. Just as significantly, because of unsolved problems with the technology of the boilers which provided steam for the engines, the performance of double-expansion engines was rarely as good in practice as it was in theory.
During the 1870s the distinction grew between 'first-class ironclads' or 'battleships' on the one hand, and 'cruising ironclads' designed for long-range work on the other. The demands on first-class ironclads for very heavy armor and armament meant increasing displacement, which reduced speed under sail; and the fashion for turrets and barbettes made a sailing rig increasingly inconvenient. , launched in 1876 but not commissioned until 1881, was the last British battleship to carry masts, and these were widely seen as a mistake. The start of the 1880s saw the end of sailing rig on ironclad battleships.
Sails persisted on 'cruising ironclads' for much longer. During the 1860s, the French navy had produced the and es as small, long-range ironclads as overseas cruisers and the British had responded with ships like of 1870. The Russian ship , laid down in 1870 and completed in 1875, was a model of a fast, long-range ironclad which was likely to be able to outrun and outfight ships like Swiftsure. Even the later , often described as the first British armored cruiser, would have been too slow to outrun General-Admiral. While Shannon was the last British ship with a retractable propeller, later armored cruisers of the 1870s retained sailing rig, sacrificing speed under steam in consequence. It took until 1881 for the Royal Navy to lay down a long-range armored warship capable of catching enemy commerce raiders, , which was completed in 1888. While sailing rigs were obsolescent for all purposes by the end of the 1880s, rigged ships were in service until the early years of the 20th century.
The final evolution of ironclad propulsion was the adoption of the triple-expansion steam engine, a further refinement which was first adopted in , laid down in 1885 and commissioned in 1891. Many ships also used a forced draught to get additional power from their engines, and this system was widely used until the introduction of the steam turbine in the middle of the first decade of the Twentieth Century.
Fleets
While ironclads spread rapidly in navies worldwide, there were few pitched naval battles involving ironclads. Most European nations settled differences on land, and the Royal Navy struggled to maintain a deterrent parity with at least France, while providing suitable protection to Britain's commerce and colonial outposts worldwide. Ironclads remained, for the British Royal Navy, a matter of defending the British Isles first and projecting power abroad second. Those naval engagements of the latter half of the 19th century which involved ironclads normally involved colonial actions or clashes between second-rate naval powers. But these encounters were often enough to convince British policy-makers of the increasing hazards of strictly naval foreign intervention, from Hampton Roads in the American Civil War to the hardening combined defences of naval arsenals such as Kronstadt and Cherbourg.
There were many types of ironclads:
Seagoing ships intended to "stand in the line of battle"; the precursors of the battleship
Coastal service and riverine vessels, including 'floating batteries' and 'monitors'
Vessels intended for commerce raiding or protection of commerce, called "armored cruisers"
Navies
The United Kingdom possessed the largest navy in the world for the whole of the ironclad period. The Royal Navy was the second to adopt ironclad warships, and it applied them worldwide in their whole range of roles. In the age of sail, the British strategy for war depended on the Royal Navy mounting a blockade of the ports of the enemy. Because of the limited endurance of steamships, this was no longer possible, so the British at times considered the risk-laden plan of engaging an enemy fleet in harbor as soon as war broke out. To this end, the Royal Navy developed a series of 'coast-defense battleships', starting with the Devastation class. These 'breastwork monitors' were markedly different from the other high-seas ironclads of the period and were an important precursor of the modern battleship. As long-range monitors they could reach Bermuda unescorted, for example. However, they were still armed with only four heavy guns and were as vulnerable to mines and obstructions (and enemy monitors) as the original monitors of the Union Navy proved to be during the Civil War.
The British prepared for an overwhelming mortar bombardment of Kronstadt by the close of the Crimean War, but never considered running the smoke-ridden, shallow-water gauntlet straight to St. Petersburg with ironclads. Likewise, monitors proved acutely unable to 'overwhelm' enemy fortifications single-handed during the American conflict, though their low-profile and heavy armor protection made them ideal for running artillery gauntlets. Mines and obstructions negated these advantages—a problem the British Admiralty frequently acknowledged but never countered throughout the period. The British never laid down enough Devastation-class 'battleships' to instantly overwhelm Cherbourg, Kronstadt or even New York City with gunfire. Although throughout the 1860s and 1870s the Royal Navy was still in many respects superior to its potential rivals, by the early 1880s widespread concern about the threat from France and Germany culminated in the Naval Defence Act, which promulgated the idea of a 'two-power standard', that Britain should possess as many ships as the next two navies combined. This standard provoked aggressive shipbuilding in the 1880s and 1890s.
British ships did not participate in any major wars in the ironclad period. The Royal Navy's ironclads only saw action as part of colonial battles or one-sided engagements like the bombardment of Alexandria in 1882. Defending British interests against Ahmed 'Urabi's Egyptian revolt, a British fleet opened fire on the fortifications around the port of Alexandria. A mixture of centre-battery and turret ships bombarded Egyptian positions for most of a day, forcing the Egyptians to retreat; return fire from Egyptian guns was heavy at first, but inflicted little damage, killing only five British sailors. Few Egyptian guns were actually dismounted, on the other hand, and the fortifications themselves were typically left intact. Had the Egyptians actually utilised the heavy mortars that were at their disposal, they might have quickly turned the tide, for the attacking British ironclads found it easy (for accuracy's sake) to simply anchor whilst firing—perfect targets for high-angle fire upon their thinly armored topdecks.
The French navy built the first ironclad to try to gain a strategic advantage over the British, but were consistently out-built by the British. Despite taking the lead with a number of innovations like breech-loading weapons and steel construction, the French navy could never match the size of the Royal Navy. In the 1870s, the construction of ironclads ceased for a while in France as the Jeune Ecole school of naval thought took prominence, suggesting that torpedo boats and unarmored cruisers would be the future of warships. Like the British, the French navy saw little action with its ironclads; the French blockade of Germany in the Franco-Prussian War was ineffective, as the war was settled entirely on land.
Russia built a number of ironclads, generally copies of British or French designs. Nonetheless, there were real innovations from Russia; the first true type of ironclad armored cruiser, General-Admiral of the 1870s, and a set of unusual but moderately-successful circular battleships referred to as "popovkas" (for Admiral Popov, who conceived the design). The Russian Navy pioneered the wide-scale use of torpedo boats during the Russo-Turkish War of 1877–1878, mainly out of necessity because of the superior numbers and quality of ironclads used by the Turkish navy. Russia expanded her navy in the 1880s and 1890s with modern armored cruisers and battleships, but the ships were manned by inexperienced crews and politically appointed leadership, which enhanced their defeat in the Battle of Tsushima on 27 May 1905.
The US Navy ended the Civil War with about fifty monitor-type coastal ironclads; by the 1870s most of these were laid up in reserve, leaving the United States virtually without an ironclad fleet. Another five large monitors were ordered in the 1870s. The limitations of the monitor type effectively prevented the US from projecting power overseas, and until the 1890s the United States would have come off badly in a conflict with even Spain or the Latin American powers. The 1890s saw the beginning of what became the Great White Fleet, and it was the modern pre-Dreadnoughts and armored cruisers built in the 1890s which defeated the Spanish fleet in the Spanish–American War of 1898. This started a new era of naval warfare.
Ironclads were widely used in South America. Both sides used ironclads in the Chincha Islands War between Spain and the combined forces of Peru and Chile in the early 1860s. The powerful Spanish participated in the Battle of Callao but was unable to inflict significant damage upon the Callao defences. Besides, Peru was able to deploy two locally built ironclads based on American Civil War designs, Loa (a wooden ship converted into a casemate ironclad) and (a small monitor armed with a single 68-pdr gun), as well as two British-built ironclads: , a centre-battery ship, and the turret ship . Numancia, was the first ironclad to circumnavigate the world under the command of Juan Bautista Antequera y Bobadilla de Eslava, arriving in Cádiz on 20 September 1867, and earning the motto: "Enloricata navis que primo terram circuivit" ["First ironclad ship to sail around the world"]). In the War of the Pacific in 1879, both Peru and Chile had ironclad warships, including some of those used a few years previously against Spain. While Independencia ran aground early on, the Peruvian ironclad made a great impact against Chilean shipping, delaying Chilean ground invasion by six months. She was eventually caught by two more modern Chilean centre-battery ironclads, and at the Battle of Angamos Point.
Ironclads were also used from the inception of the Imperial Japanese Navy (IJN). (Japanese: 甲鉄, literally "Ironclad", later renamed Azuma 東, "East") had a decisive role in the Naval Battle of Hakodate Bay in May 1869, which marked the end of the Boshin War, and the complete establishment of the Meiji Restoration. The IJN continued to develop its strength and commissioned a number of warships from British and European shipyards, first ironclads and later armored cruisers. These ships engaged the Chinese Beiyang fleet which was superior on paper at least at the Battle of the Yalu River. Thanks to superior short-range firepower, the Japanese fleet came off better, sinking or severely damaging eight ships and receiving serious damage to only four. The naval war was concluded the next year at the Battle of Weihaiwei, where the strongest remaining Chinese ships were surrendered to the Japanese.
End of the ironclad warship
There is no clearly defined end to the ironclad, besides the transition from wood hulls to all-metal. Ironclads continued to be used in World War I. Towards the end of the 19th century, the descriptions 'battleship' and 'armored cruiser' came to replace the term 'ironclad'.
The proliferation of ironclad battleship designs came to an end in the 1890s as navies reached a consensus on the design of battleships, producing the type known as the pre-dreadnought. These ships are sometimes covered in treatments of the ironclad warship. The next evolution of battleship design, the dreadnought, is never referred to as an 'ironclad'.
Legacy
H. G. Wells coined the term The Land Ironclads in a short story published in 1903, to describe fictional large armored fighting vehicles moving on pedrail wheels.
A number of ironclads have been preserved or reconstructed as museum ships.
Parts of have been recovered and are being conserved and displayed at the Mariners' Museum in Newport News, Virginia.
is today a fully restored museum ship in Portsmouth, England
is berthed at the port of Talcahuano, Chile, on display for visitors.
The is currently on display in Vicksburg, Mississippi.
Northrop Grumman in Newport News constructed a full-scale replica of . The replica was laid down in February 2005 and completed just two months later.
The Dutch Ramtorenschip (coastal ram) is currently under display in the Maritime Museum Rotterdam.
The Dutch Ramtorenschip (coastal ram) is a museum ship at Den Helder.
The complete, recovered wooden hull of , a casemate ram ironclad, is on view in Kinston, North Carolina, and, in another part of town on the Neuse River, the recreated ship, named CSS Neuse II, is nearly built and can be visited.
The hull of the casemate ironclad can be seen in the National Civil War Naval Museum in Columbus, Georgia.
A replica of the was rebuilt in 2003 as a floating museum at Weihai.
HMVS Cerberus, built 1867, has been partially sunk as a breakwater in Victoria, Australia, but is not preserved and is deteriorating in the elements.
| Technology | Naval warfare | null |
251513 | https://en.wikipedia.org/wiki/Analytic%20number%20theory | Analytic number theory | In mathematics, analytic number theory is a branch of number theory that uses methods from mathematical analysis to solve problems about the integers. It is often said to have begun with Peter Gustav Lejeune Dirichlet's 1837 introduction of Dirichlet L-functions to give the first proof of Dirichlet's theorem on arithmetic progressions. It is well known for its results on prime numbers (involving the Prime Number Theorem and Riemann zeta function) and additive number theory (such as the Goldbach conjecture and Waring's problem).
Branches of analytic number theory
Analytic number theory can be split up into two major parts, divided more by the type of problems they attempt to solve than fundamental differences in technique.
Multiplicative number theory deals with the distribution of the prime numbers, such as estimating the number of primes in an interval, and includes the prime number theorem and Dirichlet's theorem on primes in arithmetic progressions.
Additive number theory is concerned with the additive structure of the integers, such as Goldbach's conjecture that every even number greater than 2 is the sum of two primes. One of the main results in additive number theory is the solution to Waring's problem.
History
Precursors
Much of analytic number theory was inspired by the prime number theorem. Let π(x) be the prime-counting function that gives the number of primes less than or equal to x, for any real number x. For example, π(10) = 4 because there are four prime numbers (2, 3, 5 and 7) less than or equal to 10. The prime number theorem then states that x / ln(x) is a good approximation to π(x), in the sense that the limit of the quotient of the two functions π(x) and x / ln(x) as x approaches infinity is 1:
known as the asymptotic law of distribution of prime numbers.
Adrien-Marie Legendre conjectured in 1797 or 1798 that π(a) is approximated by the function a/(A ln(a) + B), where A and B are unspecified constants. In the second edition of his book on number theory (1808) he then made a more precise conjecture, with A = 1 and B ≈ −1.08366. Carl Friedrich Gauss considered the same question: "Im Jahr 1792 oder 1793" ('in the year 1792 or 1793'), according to his own recollection nearly sixty years later in a letter to Encke (1849), he wrote in his logarithm table (he was then 15 or 16) the short note "Primzahlen unter " ('prime numbers under '). But Gauss never published this conjecture. In 1838 Peter Gustav Lejeune Dirichlet came up with his own approximating function, the logarithmic integral li(x) (under the slightly different form of a series, which he communicated to Gauss). Both Legendre's and Dirichlet's formulas imply the same conjectured asymptotic equivalence of π(x) and x / ln(x) stated above, although it turned out that Dirichlet's approximation is considerably better if one considers the differences instead of quotients.
Dirichlet
Johann Peter Gustav Lejeune Dirichlet is credited with the creation of analytic number theory, a field in which he found several deep results and in proving them introduced some fundamental tools, many of which were later named after him. In 1837 he published Dirichlet's theorem on arithmetic progressions, using mathematical analysis concepts to tackle an algebraic problem and thus creating the branch of analytic number theory. In proving the theorem, he introduced the Dirichlet characters and L-functions. In 1841 he generalized his arithmetic progressions theorem from integers to the ring of Gaussian integers .
Chebyshev
In two papers from 1848 and 1850, the Russian mathematician Pafnuty L'vovich Chebyshev attempted to prove the asymptotic law of distribution of prime numbers. His work is notable for the use of the zeta function ζ(s) (for real values of the argument "s", as are works of Leonhard Euler, as early as 1737) predating Riemann's celebrated memoir of 1859, and he succeeded in proving a slightly weaker form of the asymptotic law, namely, that if the limit of π(x)/(x/ln(x)) as x goes to infinity exists at all, then it is necessarily equal to one. He was able to prove unconditionally that this ratio is bounded above and below by two explicitly given constants near to 1 for all x. Although Chebyshev's paper did not prove the Prime Number Theorem, his estimates for π(x) were strong enough for him to prove Bertrand's postulate that there exists a prime number between n and 2n for any integer n ≥ 2.
Riemann
Bernhard Riemann made some famous contributions to modern analytic number theory. In a single short paper (the only one he published on the subject of number theory), he investigated the Riemann zeta function and established its importance for understanding the distribution of prime numbers. He made a series of conjectures about properties of the zeta function, one of which is the well-known Riemann hypothesis.
Hadamard and de la Vallée-Poussin
Extending the ideas of Riemann, two proofs of the prime number theorem were obtained independently by Jacques Hadamard and Charles Jean de la Vallée-Poussin and appeared in the same year (1896). Both proofs used methods from complex analysis, establishing as a main step of the proof that the Riemann zeta function ζ(s) is non-zero for all complex values of the variable s that have the form s = 1 + it with t > 0.
Modern times
The biggest technical change after 1950 has been the development of sieve methods, particularly in multiplicative problems. These are combinatorial in nature, and quite varied. The extremal branch of combinatorial theory has in return been greatly influenced by the value placed in analytic number theory on quantitative upper and lower bounds. Another recent development is probabilistic number theory, which uses methods from probability theory to estimate the distribution of number theoretic functions, such as how many prime divisors a number has.
Specifically, the breakthroughs by Yitang Zhang, James Maynard, Terence Tao and Ben Green have all used the Goldston–Pintz–Yıldırım method, which they originally used to prove that
Developments within analytic number theory are often refinements of earlier techniques, which reduce the error terms and widen their applicability. For example, the circle method of Hardy and Littlewood was conceived as applying to power series near the unit circle in the complex plane; it is now thought of in terms of finite exponential sums (that is, on the unit circle, but with the power series truncated). The needs of Diophantine approximation are for auxiliary functions that are not generating functions—their coefficients are constructed by use of a pigeonhole principle—and involve several complex variables. The fields of Diophantine approximation and transcendence theory have expanded, to the point that the techniques have been applied to the Mordell conjecture.
Problems and results
Theorems and results within analytic number theory tend not to be exact structural results about the integers, for which algebraic and geometrical tools are more appropriate. Instead, they give approximate bounds and estimates for various number theoretical functions, as the following examples illustrate.
Multiplicative number theory
Euclid showed that there are infinitely many prime numbers. An important question is to determine the asymptotic distribution of the prime numbers; that is, a rough description of how many primes are smaller than a given number. Gauss, amongst others, after computing a large list of primes, conjectured that the number of primes less than or equal to a large number N is close to the value of the integral
In 1859 Bernhard Riemann used complex analysis and a special meromorphic function now known as the Riemann zeta function to derive an analytic expression for the number of primes less than or equal to a real number x. Remarkably, the main term in Riemann's formula was exactly the above integral, lending substantial weight to Gauss's conjecture. Riemann found that the error terms in this expression, and hence the manner in which the primes are distributed, are closely related to the complex zeros of the zeta function. Using Riemann's ideas and by getting more information on the zeros of the zeta function, Jacques Hadamard and Charles Jean de la Vallée-Poussin managed to complete the proof of Gauss's conjecture. In particular, they proved that if
then
This remarkable result is what is now known as the prime number theorem. It is a central result in analytic number theory. Loosely speaking, it states that given a large number N, the number of primes less than or equal to N is about N/log(N).
More generally, the same question can be asked about the number of primes in any arithmetic progression a + nq for any integer n. In one of the first applications of analytic techniques to number theory, Dirichlet proved that any arithmetic progression with a and q coprime contains infinitely many primes. The prime number theorem can be generalised to this problem; letting
then if a and q are coprime,
where is the totient function.
There are also many deep and wide-ranging conjectures in number theory whose proofs seem too difficult for current techniques, such as the twin prime conjecture which asks whether there are infinitely many primes p such that p + 2 is prime. On the assumption of the Elliott–Halberstam conjecture it has been proven recently that there are infinitely many primes p such that p + k is prime for some positive even k at most 12. Also, it has been proven unconditionally (i.e. not depending on unproven conjectures) that there are infinitely many primes p such that p + k is prime for some positive even k at most 246.
Additive number theory
One of the most important problems in additive number theory is Waring's problem, which asks whether it is possible, for any k ≥ 2, to write any positive integer as the sum of a bounded number of kth powers,
The case for squares, k = 2, was answered by Lagrange in 1770, who proved that every positive integer is the sum of at most four squares. The general case was proved by Hilbert in 1909, using algebraic techniques which gave no explicit bounds. An important breakthrough was the application of analytic tools to the problem by Hardy and Littlewood. These techniques are known as the circle method, and give explicit upper bounds for the function G(k), the smallest number of kth powers needed, such as Vinogradov's bound
Diophantine problems
Diophantine problems are concerned with integer solutions to polynomial equations: one may study the distribution of solutions, that is, counting solutions according to some measure of "size" or height.
An important example is the Gauss circle problem, which asks for integers points (x y) which satisfy
In geometrical terms, given a circle centered about the origin in the plane with radius r, the problem asks how many integer lattice points lie on or inside the circle. It is not hard to prove that the answer is , where as . Again, the difficult part and a great achievement of analytic number theory is obtaining specific upper bounds on the error term E(r).
It was shown by Gauss that . In general, an O(r) error term would be possible with the unit circle (or, more properly, the closed unit disk) replaced by the dilates of any bounded planar region with piecewise smooth boundary. Furthermore, replacing the unit circle by the unit square, the error term for the general problem can be as large as a linear function of r. Therefore, getting an error bound of the form
for some in the case of the circle is a significant improvement. The first to attain this was
Sierpiński in 1906, who showed . In 1915, Hardy and Landau each showed that one does not have . Since then the goal has been to show that for each fixed there exists a real number such that .
In 2000 Huxley showed that , which is the best published result.
Methods of analytic number theory
Dirichlet series
One of the most useful tools in multiplicative number theory are Dirichlet series, which are functions of a complex variable defined by an infinite series of the form
Depending on the choice of coefficients , this series may converge everywhere, nowhere, or on some half plane. In many cases, even where the series does not converge everywhere, the holomorphic function it defines may be analytically continued to a meromorphic function on the entire complex plane. The utility of functions like this in multiplicative problems can be seen in the formal identity
hence the coefficients of the product of two Dirichlet series are the multiplicative convolutions of the original coefficients. Furthermore, techniques such as partial summation and Tauberian theorems can be used to get information about the coefficients from analytic information about the Dirichlet series. Thus a common method for estimating a multiplicative function is to express it as a Dirichlet series (or a product of simpler Dirichlet series using convolution identities), examine this series as a complex function and then convert this analytic information back into information about the original function.
Riemann zeta function
Euler showed that the fundamental theorem of arithmetic implies (at least formally) the Euler product
where the product is taken over all prime numbers p.
Euler's proof of the infinity of prime numbers makes use of the divergence of the term at the left hand side for s = 1 (the so-called harmonic series), a purely analytic result. Euler was also the first to use analytical arguments for the purpose of studying properties of integers, specifically by constructing generating power series. This was the beginning of analytic number theory.
Later, Riemann considered this function for complex values of s and showed that this function can be extended to a meromorphic function on the entire plane with a simple pole at s = 1. This function is now known as the Riemann Zeta function and is denoted by ζ(s). There is a plethora of literature on this function and the function is a special case of the more general Dirichlet L-functions.
Analytic number theorists are often interested in the error of approximations such as the prime number theorem. In this case, the error is smaller than x/log x. Riemann's formula for π(x) shows that the error term in this approximation can be expressed in terms of the zeros of the zeta function. In his 1859 paper, Riemann conjectured that all the "non-trivial" zeros of ζ lie on the line but never provided a proof of this statement. This famous and long-standing conjecture is known as the Riemann Hypothesis and has many deep implications in number theory; in fact, many important theorems have been proved under the assumption that the hypothesis is true. For example, under the assumption of the Riemann Hypothesis, the error term in the prime number theorem is
In the early 20th century G. H. Hardy and Littlewood proved many results about the zeta function in an attempt to prove the Riemann Hypothesis. In fact, in 1914,
Hardy proved that there were infinitely many zeros of the zeta function on the critical line
This led to several theorems describing the density of the zeros on the critical line.
| Mathematics | Other | null |
251582 | https://en.wikipedia.org/wiki/Elephant%20shrew | Elephant shrew | Elephant shrews, also called jumping shrews or sengis, are small insectivorous mammals native to Africa, belonging to the family Macroscelididae, in the order Macroscelidea. Their traditional common English name "elephant shrew" comes from a perceived resemblance between their long noses and the trunk of an elephant, and their superficial similarity with shrews (family Soricidae) in the order Eulipotyphla. However, phylogenetic analysis has revealed that elephant shrews are not properly classified with true shrews, but are in fact more closely related to elephants than to shrews. In 1997, the biologist Jonathan Kingdon proposed that they instead be called "sengis" (singular sengi), a term derived from the Bantu languages of Africa, and in 1998, they were classified into the new clade Afrotheria.
They are widely distributed across the southern part of Africa, and although common nowhere, can be found in almost any type of habitat, from the Namib Desert to boulder-strewn outcrops in South Africa to thick forest. One species, the North African elephant shrew, remains in the semi-arid, mountainous country in the far northwest of Africa. The Somali elephant shrew went unobserved from 1968 to 2020 but was rediscovered by a group of scientists in Djibouti.
Description
Elephant shrews are small, quadrupedal, insectivorous mammals. They have scaly tails, long snouts, and bear a superficial resemblance to shrews or rodents. They have long legs relative to their size, which are used to move from one place to another like rabbits.
Elephant shrews use their flexible proboscises to search for food, with the length of the snout varying between species.
They are one of the fastest small mammals, having been recorded to reach speeds of . They vary in size from about , from . One species of giant sengi, the grey-faced sengi, weighs about 700 g. Compared to other mammalian insectivores, sengis have relatively large brains. Their lifespans are about two and a half to four years in the wild. They have large canine teeth, and also high-crowned cheek teeth similar to those of ungulates. Their dental formula is G
Behavior and ecology
Although mostly diurnal and very active, they are difficult to trap and very seldom seen; elephant shrews are wary, well camouflaged, and adept at dashing away from threats. Several species make a series of cleared pathways through the undergrowth and spend their day patrolling them for insect life. If the animal is disturbed, the pathway provides an obstacle-free escape route.
Elephant shrews are solitary animals, despite many species living in monogamous pairs. They share and defend their home territory, which is marked using their scent glands. Scent markings are also used for mate attraction.
Short-eared elephant shrews inhabit the dry steppes and stone deserts of southwestern Africa. They can even be found in the Namib Desert, one of the driest regions of the earth. Females drive away other females, while males try to ward off other males. Although they live in pairs, the partners do not care much for each other and their sole purpose of even associating with the opposite sex is for reproduction. Social behaviors are not very common and they even have separate nests. The one or two young are well developed at birth; they are able to run within a few hours.
Female elephant shrews undergo a menstrual cycle similar to that of human females, making it one of the few nonprimate mammals to do so. Elephant shrews were used in the 1940s to study the human menstruation cycle. The elephant shrew mating period lasts for several days. After mating, the pair will return to their solitary habits. After a gestation period varying from 45 to 60 days, the female will bear litters of one to three young several times a year. The young are born relatively well developed, but remain in the nest for several days before venturing outside.
After five days, the young's milk diet is supplemented with mashed insects, which are collected and transported in the cheek pouches of the female. The young then slowly start to explore their environment and hunt for insects. After about 15 days, the young will begin the migratory phase of their lives, which lessens their dependency on their mother. The young will then establish their own home ranges (about ) and will become sexually active within 41–46 days.
The thermal characteristics of elephant shrews with similar body size, habitat and distribution are very close in most of the classifications. They can maintain homeothermy in different ambient temperatures where most of the species regulate their body temperature at 35 °C and neither become hyperthermic but they balance the heat offload by increasing the EWL (evaporative water loss).
Feeding habits
Elephant shrews mainly eat insects, spiders, centipedes, millipedes, and earthworms. While awake, as much as 80% of their time may be spent foraging. An elephant shrew uses its nose to find prey and uses its tongue to flick small food into its mouth, much like an anteater. Eating large prey can pose a challenge; an elephant shrew struggling with an earthworm must first pin its prey to the ground with a forefoot. Then, turning its head to one side, it chews pieces off with its cheek teeth, much like a dog chewing a bone. This is a sloppy process, and many small pieces of worm drop to the ground; these are simply flicked up with the tongue. Some elephant shrews also feed on small amounts of plant matter, especially new leaves, seeds, and small fruits.
Evolution
A number of fossil species are known, all from Africa. They were separate from the similar-appearing order Leptictida. A considerable diversification of macroscelids occurred in the Paleogene period. The earliest definitive member was Chambia from the early-middle Eocene of Tunisia. Some early macroscelids, such as Myohyrax, were so similar to hyraxes that they were initially included with that group, while others, such as Mylomygale, were relatively rodent-like. These unusual forms all died out by the Pleistocene. Although macroscelids were classified in the past with many groups, often on the basis of superficial characteristics, considerable morphological and molecular evidence places them within Afrotheria, at the base of Afroinsectivora.
In terms of timing, the divergence between macroscelids and afrosoricidans is thought to have occurred roughly 57.5 million years (Ma) ago, in the late Paleocene, while the diversification of extant macroscelids apparently began when the Rhynchocyon lineage split off about 33 Ma ago, in the early Oligocene. Elephantulus is considered to have separated from Macroscelidini later in the Oligocene, about 28.5 Ma ago.
Phylogeny
Classification
The 20 species of elephant shrew are placed in six genera, three of which are monotypic:
ORDER MACROSCELIDEA
Family Macroscelididae
Genus Elephantulus
Short-snouted elephant shrew, E. brachyrhynchus
Cape elephant shrew, E. edwardii
Dusky-footed elephant shrew, E. fuscipes
Dusky elephant shrew, E. fuscus
Bushveld elephant shrew, E. intufi
Eastern rock elephant shrew, E. myurus
Karoo rock elephant shrew, E. pilicaudus
Western rock elephant shrew, E. rupestris
Genus Galegeeska
Somali elephant shrew, G. revoilii
Rufous elephant shrew, G. rufescens
Genus Macroscelides
Namib round-eared sengi, M. flavicaudatus
Etendeka round-eared sengi, M. micus
Round-eared elephant shrew, M. proboscideus
Genus Petrodromus
Four-toed elephant shrew, P. tetradactylus
Genus Petrosaltator
North African elephant shrew, P. rozeti
Genus Rhynchocyon
Golden-rumped sengi, R. chrysopygus
Rhynchocyon chrysopygus mandelai
Chequered sengi, R. cirnei
Rhynchocyon cirnei cirnei
Rhynchocyon cirnei shirensis
Rhynchocyon cirnei reichardi
Rhynchocyon cirnei hendersoni
Rhynchocyon cirnei macrurus
Black and rufous sengi, R. petersi
Rhynchocyon petersi petersi
Rhynchocyon petersi adersi
Stuhlmann's sengi, R. stuhlmanni
Grey-faced sengi, R. udzungwensis
| Biology and health sciences | Other afrotheres | Animals |
251647 | https://en.wikipedia.org/wiki/Nut%20%28fruit%29 | Nut (fruit) | A nut is a fruit consisting of a hard or tough nutshell protecting a kernel which is usually edible. In general usage and in a culinary sense, many dry seeds are called nuts, but in a botanical context, "nut" implies that the shell does not open to release the seed (indehiscent).
Most seeds come from fruits that naturally free themselves from the shell, but this is not the case in nuts such as hazelnuts, chestnuts, and acorns, which have hard shell walls and originate from a compound ovary.
Definition
A seed is the mature fertilised ovule of a plant; it consists of three parts, the embryo which will develop into a new plant, stored food for the embryo, and a protective seed coat. Botanically, a nut is a fruit with a woody pericarp developing from a syncarpous gynoecium. Nuts may be contained in an involucre, a cup-shaped structure formed from the flower bracts. The involucre may be scaly, spiny, leafy or tubular, depending on the species of nut. Most nuts come from the pistils with inferior ovaries (see flower) and all are indehiscent (not opening at maturity). True nuts are produced, for example, by some plant families of the order Fagales. These include beech (Fagus), chestnut (Castanea), oak (Quercus), stone-oak (Lithocarpus) and tanoak (Notholithocarpus) in the family Fagaceae, as well as hazel, filbert (Corylus) and hornbeam (Carpinus) in the family Betulaceae.
A small nut may be called a "nutlet" (formerly called a nucule, a term otherwise referring to the oogonium of stoneworts). In botany, the term "nutlet" can be used to describe a pyrena or pyrene, which is a seed covered by a stony layer, such as the kernel of a drupe.
Walnuts and hickories (Juglandaceae) have fruits that are difficult to classify. They are considered to be nuts under some definitions but are also referred to as drupaceous nuts.
Evolutionary history
Toxicity
Nuts used for food are a common source of food allergens. Reactions can range from mild symptoms to severe ones, a condition known as anaphylaxis, which can be life-threatening. The reaction is due to the release of histamine by the body in response to an allergen in the nuts, causing skin and other possible reactions. Many experts suggest that a person with an allergy to peanuts should avoid eating tree nuts, and vice versa.
Consumption as food
Nuts contain the diverse nutrients that are needed for the growth of a new plant. Composition varies, but they tend to have a low water and carbohydrate content, with high levels of fats, protein, dietary minerals, and vitamins.
Nuts are eaten by humans and wildlife. Because nuts generally have a high oil content, they are a significant energy source. Many seeds are edible by humans and used in cooking, eaten raw, sprouted, or roasted as a snack food, ground to make nut butters, or pressed for oil that is used in cooking and cosmetics.
Constituents
Nuts are the source of energy and nutrients for the new plant. They contain a relatively large quantity of calories, essential unsaturated and monounsaturated fats including linoleic acid and linolenic acid, vitamins, and essential amino acids.
| Biology and health sciences | Plant: General | null |
251662 | https://en.wikipedia.org/wiki/Golden%20mole | Golden mole | Golden moles are small insectivorous burrowing mammals endemic to Sub-Saharan Africa. They comprise the family Chrysochloridae (the only family in the suborder Chrysochloridea) and as such they are taxonomically distinct from the true moles, family Talpidae, and other mole-like families, all of which, to various degrees, they resemble as a result of evolutionary convergence. There are 21 species. Some (e.g., Chrysochloris asiatica, Amblysomus hottentotus) are relatively common, whereas others (e.g., species of Chrysospalax, Cryptochloris, Neamblysomus) are rare and endangered.
Taxonomy
Despite their extreme evolutionary convergence with moles, their closest relatives are the otter shrews and tenrecs. They are more distantly related to other insectivorous African mammals such as aardvarks and elephant shrews, and even more distantly related to elephants and sirenians.
Characteristics and affinities
Like most burrowing mammals with similar habits, the Chrysochloridae have short legs with powerful digging claws, very dense fur that repels dirt and moisture, and toughened skin, particularly on the head. The fur varies from black to pale yellow or grey, and has an iridescent sheen of green, blue, violet or copper when viewed in light.
Their eyes are non-functional and covered with furred skin. The external ears are just tiny openings. In particular, golden moles bear a remarkable resemblance to the marsupial moles of Australia, family Notoryctidae, which they resemble so suggestively that at one time, the marsupial/placental divide notwithstanding, some argued that they were related. Considerations that influenced the debate might have included the view that the Chrysochloridae are very primitive placentals and the fact that they have many mole-like specializations similar to specializations in marsupial moles.
The rhinarium is a greatly enlarged, dry leathery pad that protects their nostrils while the animal digs. In this respect, too, they resemble the marsupial moles. Some authors claim their primary sense is of touch, and they are particularly sensitive to vibrations, which may indicate approaching danger. Note below, however, the observations on the malleus in the middle ear.
The species range in size from about to about . They have muscular shoulders and the forelimbs are radically adapted for digging; all the toes on the forefeet have been reduced, except for a large, pick-like third claw on the third toe. In comparison to true moles, the fifth digit is absent and the first and fourth digits are vestigial. The adaptations of the hind feet are less dramatic: They retain all five toes and are webbed as an adaptation to efficient backward shoveling of soil loosened by the front claws.
At one time, the Chrysochloridae were regarded as primitive. Supporting arguments of this included that they were thought to have originated in Gondwana, that they had a low resting metabolic rate, and that they could switch off thermoregulation when inactive. Like the tenrecs, they possess a cloaca, and males lack a scrotum. However, these points are no longer regarded as strongly suggestive of golden moles as undeveloped "reptilian mammals"; some are seen rather as adaptations to regional climatic conditions. Going into a torpor when resting or during cold weather enables them to conserve energy and reduce urgent requirements for food. Similarly, they have developed particularly efficient kidneys, and some species do not need to drink water at all. They tend to drown if they fall into water.
Habits and ecology
Chrysochloridae are subterranean, afrotherian mammals endemic to sub-Saharan Africa, and most of which are recorded from South Africa in particular. Other regions include Lake Victoria, Western Cape, and Namibia. They live in a variety of environments; forest, swamps, deserts, or mountainous terrain. Chrysospalax species tend to forage above ground in leaf litter in forests or in meadows. Eremitalpa species such as Grant's golden mole live in the sandy Namib desert, where they cannot form tunnels because the sand collapses. Instead during the day, when they must seek shelter, they "swim" through the loose sand, using their broad claws to paddle, and dive down some to where it is bearably cool. There they enter a state of torpor, thus conserving energy. At night they emerge to forage on the surface rather than wasting energy shifting sand. Their main prey are termites that live under isolated grass clumps, and they might travel for a night in search of food. They seek promising clumps by listening for wind-rustled grass-root stresses and termites' head-banging alarm signals, neither of which can be heard easily above ground, so they stop periodically and dip their heads under the sand to listen.
Most other species construct both foraging superficial burrows and deeper permanent burrows for residence. Residential burrows are relatively complex in form and may penetrate as far as below ground and include deep chambers for use for refuge, and other chambers as latrines. They push excavated soil up to the surface, as in mole-hills, or compact it into the tunnel walls.
They feed on small insects and earthworms or small vertebrates such as lizards or burrowing snakes. They depend on their sense of hearing to locate much of their prey, and the cochleas of a number of golden mole species have been found to be long and highly coiled, which may indicate a greater ecological dependence on low frequency auditory cues than seen in Talpid moles.
Morphology
Golden moles share a number of features, varying by species, seldom seen elsewhere among living mammals, including three forearm long-bones, hyoid-mandible articulation, and a hypertrophied malleus. Some species have hypertrophied (enlarged) middle ear ossicles, in particular the malleus. These animals have the largest malleus relative to body size of any animal. This morphology may be adapted for the detection of seismic signals. In this respect there is some apparent convergent evolution to burrowing reptiles in the family Amphisbaenidae.
Reproduction
Females give birth to one to three hairless young in a grass-lined nest within the burrow system. Breeding occurs throughout the year. The adults are solitary, and their burrowing territory may be aggressively defended from intruders, especially where resources are relatively scarce.
Status
Of the 21 species of golden mole, no fewer than 11 are threatened with extinction, the primary cause of which being human-induced habitat loss. Additionally sand mining, poor agricultural practices, and predation by domestic cats and dogs are causes of population decline.
Classification
The taxonomy of the Chrysochloridae is undergoing a review in the light of new genetic information. They have traditionally been listed with the shrews, hedgehogs and a grab-bag of small, difficult-to-place creatures as part of the order Insectivora. Some authorities retain this classification, at least for the time being. Others group the golden moles with the tenrecs in a new order, which is sometimes known as Tenrecomorpha, while others call it Afrosoricida and reserve Tenrecomorpha for the family Tenrecidae.
ORDER AFROSORICIDA
Suborder Tenrecomorpha
Family Tenrecidae: tenrecs, 34 species in 10 genera
Suborder Chrysochloridea
Family Chrysochloridae
Subfamily Chrysochlorinae
Genus Carpitalpa
Arends's golden mole (Carpitalpa arendsi)
Genus Chlorotalpa
Duthie's golden mole (Chlorotalpa duthieae)
Sclater's golden mole (Chlorotalpa sclateri)
Genus Chrysochloris
Subgenus Chrysochloris
Cape golden mole (Chrysochloris asiatica)
Visagie's golden mole (Chrysochloris visagiei)
Subgenus Kilimatalpa
Stuhlmann's golden mole (Chrysochloris stuhlmanni)
Genus Chrysospalax
Giant golden mole (Chrysospalax trevelyani)
Rough-haired golden mole (Chrysospalax villosus)
Genus Cryptochloris
De Winton's golden mole (Cryptochloris wintoni)
Van Zyl's golden mole (Cryptochloris zyli)
Genus Eremitalpa
Grant's golden mole (Eremitalpa granti)
Subfamily Amblysominae
Genus Amblysomus
Fynbos golden mole (Amblysomus corriae)
Hottentot golden mole (Amblysomus hottentotus)
Marley's golden mole (Amblysomus marleyi)
Robust golden mole (Amblysomus robustus)
Highveld golden mole (Amblysomus septentrionalis)
Genus Calcochloris
Subgenus Calcochloris
Yellow golden mole (Calcochloris obtusirostris)
Subgenus incertae sedis
Somali golden mole (Calcochloris tytonis)
Genus Huetia
Congo golden mole (Huetia leucorhina)
Genus Neamblysomus
Juliana's golden mole (Neamblysomus julianae)
Gunning's golden mole (Neamblysomus gunningi)
| Biology and health sciences | Other afrotheres | Animals |
251720 | https://en.wikipedia.org/wiki/Praseodymium | Praseodymium | Praseodymium is a chemical element; it has symbol Pr and the atomic number 59. It is the third member of the lanthanide series and is considered one of the rare-earth metals. It is a soft, silvery, malleable and ductile metal, valued for its magnetic, electrical, chemical, and optical properties. It is too reactive to be found in native form, and pure praseodymium metal slowly develops a green oxide coating when exposed to air.
Praseodymium always occurs naturally together with the other rare-earth metals. It is the sixth-most abundant rare-earth element and fourth-most abundant lanthanide, making up 9.1 parts per million of the Earth's crust, an abundance similar to that of boron. In 1841, Swedish chemist Carl Gustav Mosander extracted a rare-earth oxide residue he called didymium from a residue he called "lanthana", in turn separated from cerium salts. In 1885, the Austrian chemist Carl Auer von Welsbach separated didymium into two elements that gave salts of different colours, which he named praseodymium and neodymium. The name praseodymium comes from the Ancient Greek (), meaning 'leek-green', and () 'twin'.
Like most rare-earth elements, praseodymium most readily forms the +3 oxidation state, which is the only stable state in aqueous solution, although the +4 oxidation state is known in some solid compounds and, uniquely among the lanthanides, the +5 oxidation state is attainable in matrix-isolation conditions. The 0, +1, and +2 oxidation states are rarely found. Aqueous praseodymium ions are yellowish-green, and similarly, praseodymium results in various shades of yellow-green when incorporated into glasses. Many of praseodymium's industrial uses involve its ability to filter yellow light from light sources.
Physical properties
Praseodymium is the third member of the lanthanide series, and a member of the rare-earth metals. In the periodic table, it appears between the lanthanides cerium to its left and neodymium to its right, and above the actinide protactinium. It is a ductile metal with a hardness comparable to that of silver. Praseodymium is calculated to have a very large atomic radius; with a radius of 247 pm, barium, rubidium and caesium are larger. However, observationally, it is usually 185 pm.
Neutral praseodymium's 59 electrons are arranged in the configuration [Xe]4f36s2.
Like most other lanthanides, praseodymium usually uses only three electrons as valence electrons, as the remaining 4f electrons are too strongly bound to engage in bonding: this is because the 4f orbitals penetrate the most through the inert xenon core of electrons to the nucleus, followed by 5d and 6s, and this penetration increases with higher ionic charge. Even so, praseodymium can in some compounds lose a fourth valence electron because it is early in the lanthanide series, where the nuclear charge is still low enough and the 4f subshell energy high enough to allow the removal of further valence electrons.
Similarly to the other early lanthanides, praseodymium has a double hexagonal close-packed crystal structure at room temperature, called the alpha phase (α-Pr). At it transforms to a different allotrope that has a body-centered cubic structure (β-Pr), and it melts at .
Praseodymium, like all of the lanthanides, is paramagnetic at room temperature. Unlike some other rare-earth metals, which show antiferromagnetic or ferromagnetic ordering at low temperatures, praseodymium is paramagnetic at all temperatures above 1 K.
Chemical properties
Praseodymium metal tarnishes slowly in air, forming a spalling green oxide layer like iron rust; a centimetre-sized sample of praseodymium metal corrodes completely in about a year. It burns readily at 150 °C to form praseodymium(III,IV) oxide, a nonstoichiometric compound approximating to Pr6O11:
12 Pr + 11 O2 → 2 Pr6O11
This may be reduced to praseodymium(III) oxide (Pr2O3) with hydrogen gas. Praseodymium(IV) oxide, PrO2, is the most oxidised product of the combustion of praseodymium and can be obtained by either reaction of praseodymium metal with pure oxygen at 400 °C and 282 bar or by disproportionation of Pr6O11 in boiling acetic acid. The reactivity of praseodymium conforms to periodic trends, as it is one of the first and thus one of the largest lanthanides. At 1000 °C, many praseodymium oxides with composition PrO2−x exist as disordered, nonstoichiometric phases with 0 < x < 0.25, but at 400–700 °C the oxide defects are instead ordered, creating phases of the general formula PrnO2n−2 with n = 4, 7, 9, 10, 11, 12, and ∞. These phases PrOy are sometimes labelled α and β′ (nonstoichiometric), β (y = 1.833), δ (1.818), ε (1.8), ζ (1.778), ι (1.714), θ, and σ.
Praseodymium is an electropositive element and reacts slowly with cold water and quite quickly with hot water to form praseodymium(III) hydroxide:
2 Pr (s) + 6 H2O (l) → 2 Pr(OH)3 (aq) + 3 H2 (g)
Praseodymium metal reacts with all the stable halogens to form trihalides:
2 Pr (s) + 3 F2 (g) → 2 PrF3 (s) [green]
2 Pr (s) + 3 Cl2 (g) → 2 PrCl3 (s) [green]
2 Pr (s) + 3 Br2 (g) → 2 PrBr3 (s) [green]
2 Pr (s) + 3 I2 (g) → 2 PrI3 (s)
The tetrafluoride, PrF4, is also known, and is produced by reacting a mixture of sodium fluoride and praseodymium(III) fluoride with fluorine gas, producing Na2PrF6, following which sodium fluoride is removed from the reaction mixture with liquid hydrogen fluoride. Additionally, praseodymium forms a bronze diiodide; like the diiodides of lanthanum, cerium, and gadolinium, it is a praseodymium(III) electride compound.
Praseodymium dissolves readily in dilute sulfuric acid to form solutions containing the chartreuse Pr3+ ions, which exist as [Pr(H2O)9]3+ complexes:
2 Pr (s) + 3 H2SO4 (aq) → 2 Pr3+ (aq) + 3 (aq) + 3 H2 (g)
Dissolving praseodymium(IV) compounds in water does not result in solutions containing the yellow Pr4+ ions; because of the high positive standard reduction potential of the Pr4+/Pr3+ couple at +3.2 V, these ions are unstable in aqueous solution, oxidising water and being reduced to Pr3+. The value for the Pr3+/Pr couple is −2.35 V. However, in highly basic aqueous media, Pr4+ ions can be generated by oxidation with ozone.
Although praseodymium(V) in the bulk state is unknown, the existence of praseodymium in its +5 oxidation state (with the stable electron configuration of the preceding noble gas xenon) under noble-gas matrix isolation conditions was reported in 2016. The species assigned to the +5 state were identified as [PrO2]+, its O2 and Ar adducts, and PrO2(η2-O2).
Organopraseodymium compounds
Organopraseodymium compounds are very similar to those of the other lanthanides, as they all share an inability to undergo π backbonding. They are thus mostly restricted to the mostly ionic cyclopentadienides (isostructural with those of lanthanum) and the σ-bonded simple alkyls and aryls, some of which may be polymeric. The coordination chemistry of praseodymium is largely that of the large, electropositive Pr3+ ion, and is thus largely similar to those of the other early lanthanides La3+, Ce3+, and Nd3+. For instance, like lanthanum, cerium, and neodymium, praseodymium nitrates form both 4:3 and 1:1 complexes with 18-crown-6, whereas the middle lanthanides from promethium to gadolinium can only form the 4:3 complex and the later lanthanides from terbium to lutetium cannot successfully coordinate to all the ligands. Such praseodymium complexes have high but uncertain coordination numbers and poorly defined stereochemistry, with exceptions resulting from exceptionally bulky ligands such as the tricoordinate [Pr{N(SiMe3)2}3]. There are also a few mixed oxides and fluorides involving praseodymium(IV), but it does not have an appreciable coordination chemistry in this oxidation state like its neighbour cerium. However, the first example of a molecular complex of praseodymium(IV) has recently been reported.
Isotopes
Praseodymium has only one stable and naturally occurring isotope, 141Pr. It is thus a mononuclidic and monoisotopic element, and its standard atomic weight can be determined with high precision as it is a constant of nature. This isotope has 82 neutrons, which is a magic number that confers additional stability. This isotope is produced in stars through the s- and r-processes (slow and rapid neutron capture, respectively). Thirty-eight other radioisotopes have been synthesized. All of these isotopes have half-lives under a day (and most under a minute), with the single exception of 143Pr with a half-life of 13.6 days. Both 143Pr and 141Pr occur as fission products of uranium. The primary decay mode of isotopes lighter than 141Pr is positron emission or electron capture to isotopes of cerium, while that of heavier isotopes is beta decay to isotopes of neodymium.
History
In 1751, the Swedish mineralogist Axel Fredrik Cronstedt discovered a heavy mineral from the mine at Bastnäs, later named cerite. Thirty years later, the fifteen-year-old Wilhelm Hisinger, from the family owning the mine, sent a sample of it to Carl Scheele, who did not find any new elements within. In 1803, after Hisinger had become an ironmaster, he returned to the mineral with Jöns Jacob Berzelius and isolated a new oxide, which they named ceria after the dwarf planet Ceres, which had been discovered two years earlier. Ceria was simultaneously and independently isolated in Germany by Martin Heinrich Klaproth. Between 1839 and 1843, ceria was shown to be a mixture of oxides by the Swedish surgeon and chemist Carl Gustaf Mosander, who lived in the same house as Berzelius; he separated out two other oxides, which he named lanthana and didymia. He partially decomposed a sample of cerium nitrate by roasting it in air and then treating the resulting oxide with dilute nitric acid. The metals that formed these oxides were thus named lanthanum and didymium.
While lanthanum turned out to be a pure element, didymium was not and turned out to be only a mixture of all the stable early lanthanides from praseodymium to europium, as had been suspected by Marc Delafontaine after spectroscopic analysis, though he lacked the time to pursue its separation into its constituents. The heavy pair of samarium and europium were only removed in 1879 by Paul-Émile Lecoq de Boisbaudran and it was not until 1885 that Carl Auer von Welsbach separated didymium into praseodymium and neodymium. Von Welsbach confirmed the separation by spectroscopic analysis, but the products were of relatively low purity. Since neodymium was a larger constituent of didymium than praseodymium, it kept the old name with disambiguation, while praseodymium was distinguished by the leek-green colour of its salts (Greek πρασιος, "leek green"). The composite nature of didymium had previously been suggested in 1882 by Bohuslav Brauner, who did not experimentally pursue its separation.
Occurrence and production
Praseodymium is not particularly rare, despite it being in the rare-earth metals, making up 9.2 mg/kg of the Earth's crust. Praseodymium's classification as a rare-earth metal comes from its rarity relative to "common earths" such as lime and magnesia, the few known minerals containing it for which extraction is commercially viable, as well as the length and complexity of extraction. Although not particularly rare, praseodymium is never found as a dominant rare earth in praseodymium-bearing minerals. It is always preceded by cerium and lanthanum and usually also by neodymium.
The Pr3+ ion is similar in size to the early lanthanides of the cerium group (those from lanthanum up to samarium and europium) that immediately follow in the periodic table, and hence it tends to occur along with them in phosphate, silicate and carbonate minerals, such as monazite (MIIIPO4) and bastnäsite (MIIICO3F), where M refers to all the rare-earth metals except scandium and the radioactive promethium (mostly Ce, La, and Y, with somewhat less Nd and Pr). Bastnäsite is usually lacking in thorium and the heavy lanthanides, and the purification of the light lanthanides from it is less involved. The ore, after being crushed and ground, is first treated with hot concentrated sulfuric acid, evolving carbon dioxide, hydrogen fluoride, and silicon tetrafluoride. The product is then dried and leached with water, leaving the early lanthanide ions, including lanthanum, in solution.
The procedure for monazite, which usually contains all the rare earth, as well as thorium, is more involved. Monazite, because of its magnetic properties, can be separated by repeated electromagnetic separation. After separation, it is treated with hot concentrated sulfuric acid to produce water-soluble sulfates of rare earth. The acidic filtrates are partially neutralized with sodium hydroxide to pH 3–4, during which thorium precipitates as hydroxide and is removed. The solution is treated with ammonium oxalate to convert rare earth to their insoluble oxalates, the oxalates are converted to oxides by annealing, and the oxides are dissolved in nitric acid. This last step excludes one of the main components, cerium, whose oxide is insoluble in HNO3. Care must be taken when handling some of the residues as they contain 228Ra, the daughter of 232Th, which is a strong gamma emitter.
Praseodymium may then be separated from the other lanthanides via ion-exchange chromatography, or by using a solvent such as tributyl phosphate where the solubility of Ln3+ increases as the atomic number increases. If ion-exchange chromatography is used, the mixture of lanthanides is loaded into one column of cation-exchange resin and Cu2+ or Zn2+ or Fe3+ is loaded into the other. An aqueous solution of a complexing agent, known as the eluant (usually triammonium edtate), is passed through the columns, and Ln3+ is displaced from the first column and redeposited in a compact band at the top of the column before being re-displaced by . The Gibbs free energy of formation for Ln(edta·H) complexes increases along with the lanthanides by about one quarter from Ce3+ to Lu3+, so that the Ln3+ cations descend the development column in a band and are fractionated repeatedly, eluting from heaviest to lightest. They are then precipitated as their insoluble oxalates, burned to form the oxides, and then reduced to metals.
Applications
Leo Moser (not to be confused with the mathematician of the same name), son of Ludwig Moser, founder of the Moser Glassworks in what is now Karlovy Vary in the Czech Republic, investigated the use of praseodymium in glass coloration in the late 1920s, yielding a yellow-green glass given the name "Prasemit". However, at that time far cheaper colorants could give a similar color, so Prasemit was not popular, few pieces were made, and examples are now extremely rare. Moser also blended praseodymium with neodymium to produce "Heliolite" glass ("Heliolit" in German), which was more widely accepted. The first enduring commercial use of purified praseodymium, which continues today, is in the form of a yellow-orange "Praseodymium Yellow" stain for ceramics, which is a solid solution in the zircon lattice. This stain has no hint of green in it; by contrast, at sufficiently high loadings, praseodymium glass is distinctly green rather than pure yellow.
Like many other lanthanides, praseodymium's shielded f-orbitals allow for long excited state lifetimes and high luminescence yields. Pr3+ as a dopant ion therefore sees many applications in optics and photonics. These include DPSS-lasers, single-mode fiber optical amplifiers, fiber lasers, upconverting nanoparticles as well as activators in red, green, blue, and ultraviolet phosphors. Silicate crystals doped with praseodymium ions have also been used to slow a light pulse down to a few hundred meters per second.
As the lanthanides are so similar, praseodymium can substitute for most other lanthanides without significant loss of function, and indeed many applications such as mischmetal and ferrocerium alloys involve variable mixes of several lanthanides, including small quantities of praseodymium. The following more modern applications involve praseodymium specifically or at least praseodymium in a small subset of the lanthanides:
In combination with neodymium, another rare-earth element, praseodymium is used to create high-power magnets notable for their strength and durability. In general, most alloys of the cerium-group rare earths (lanthanum through samarium) with 3d transition metals give extremely stable magnets that are often used in small equipment, such as motors, printers, watches, headphones, loudspeakers, and magnetic storage.
Praseodymium–nickel intermetallic (PrNi5) has such a strong magnetocaloric effect that it has allowed scientists to approach within one thousandth of a degree of absolute zero.
As an alloying agent with magnesium to create high-strength metals that are used in aircraft engines; yttrium and neodymium are suitable substitutes.
Praseodymium is present in the rare-earth mixture whose fluoride forms the core of carbon arc lights, which are used in the motion picture industry for studio lighting and projector lights.
Praseodymium compounds give glasses, enamels and ceramics a yellow color.
Praseodymium is a component of didymium glass, which is used to make certain types of welder's and glass blower's goggles.
Praseodymium oxide in solid solution with ceria or ceria-zirconia has been used as an oxidation catalyst.
Due to its role in permanent magnets used for wind turbines, it has been argued that praseodymium will be one of the main objects of geopolitical competition in a world running on renewable energy. However, this perspective has been criticized for failing to recognize that most wind turbines do not use permanent magnets and for underestimating the power of economic incentives for expanded production.
Biological role and precautions
The early lanthanides have been found to be essential to some methanotrophic bacteria living in volcanic mudpots, such as Methylacidiphilum fumariolicum: lanthanum, cerium, praseodymium, and neodymium are about equally effective. Praseodymium is otherwise not known to have a biological role in any other organisms, but it is not very toxic either. Intravenous injection of rare earths into animals has been known to impair liver function, but the main side effects from inhalation of rare-earth oxides in humans come from radioactive thorium and uranium impurities.
| Physical sciences | Chemical elements_2 | null |
251882 | https://en.wikipedia.org/wiki/Quart | Quart | The quart (symbol: qt) is a unit of volume equal to a quarter of a gallon. Three kinds of quarts are currently used: the liquid quart and dry quart of the US customary system and the of the British imperial system. All are roughly equal to one liter. It is divided into two pints or (in the US) four cups. Historically, the exact size of the quart has varied with the different values of gallons over time and in reference to different commodities.
Name
The term comes from the Latin (meaning one-quarter) via the French . However, although the French word has the same root, it frequently means something entirely different. In Canadian French in particular, the quart is called , whilst the pint is called .
History
Since gallons of various sizes have historically been in use, the corresponding quarts have also existed with various sizes.
Definitions and equivalencies
US liquid quart
In the United States, traditional length and volume measures have been legally standardized for commerce by the international yard and pound agreement of 1959, using the definition of 1 yard being exactly equal to 0.9144 meters. From this definition is derived the metric equivalencies for inches, feet, and miles, area measures, and measures of volume. The US liquid quart equals 57.75 cubic inches, which is exactly equal to .
US dry quart
In the United States, the dry quart is equal to one quarter of a US dry gallon, or exactly .
Imperial quart
The imperial quart, which is used for both liquid and dry capacity, is equal to one quarter of an imperial gallon, or exactly 1.1365225 liters. In the United Kingdom goods may be sold by the quart if the equivalent metric measure is also given.
In Canadian French, by federal law, the imperial quart is called .
Winchester quart
The Winchester quart is an archaic measure, roughly equal to 2 imperial quarts or 2.25 liters. The 2.5L bottles in which laboratory chemicals are supplied are sometimes referred to as Winchester quart bottles, although they contain slightly more than a traditional Winchester quart.
Reputed quart
The reputed quart was a measure equal to two-thirds of an imperial quart (or one-sixth of an imperial gallon), at about 0.7577liters, which is very close to one US fifth (0.757 liters).
The reputed quart was previously recognized as a standard size of wine bottle in the United Kingdom, and is only about 1% larger than the current standard wine bottle of 0.75L.
| Physical sciences | Volume | Basics and measurement |
251887 | https://en.wikipedia.org/wiki/Luxury%20car | Luxury car | A luxury car is a passenger vehicle providing above-average to high-end comfort levels, features, and equipment. More expensive materials and surface finishes are often used, and buyers expect better build quality. The usually higher pricing and more upscale appearance are often associated with the users' higher social status compared to low- and mid-market segment cars.
The term is relative and partially subjective, reflecting both objective qualities of a car, as well as the projected and perceived image of the vehicle's marque. Luxury brands rank above premium brands, though there is no clear distinction between the two.
Most luxury cars were large, though smaller sports-oriented models were produced. "Compact" luxury vehicles such as hatchbacks and off-road capable sport utility vehicles are more recent expansions of luxury qualities in various cars.
Classification standards
Several car classification schemes include a luxury category, such as:
Australia: Since the year 2000, the Federal Government's luxury car tax applies to new vehicles over a certain purchase price, with higher thresholds applying for cars considered as fuel efficient. As of 2019, the thresholds were approximately AU$66,000 (US$,000) for normal cars and AU$76,000 (US$,000) for fuel-efficient cars.
Europe: Luxury cars are classified as F-segment vehicles in the European Commission classification scheme.
Italy: The term "auto di lusso" is used for luxury cars.
France: The term "voiture de luxe" is used for luxury cars.
Germany: The term (upper class) is used for luxury cars.
Russia: The term ( ("representative class vehicle, also translated as luxury vehicle) is used for luxury cars.
Rental cars: The ACRISS Car Classification Code is a system used by many car rental companies to define equivalent vehicles across brands. This system includes "Luxury" and "Luxury Elite" categories (along with "Premium" and "Premium Elite" categories). The criteria for a vehicle to be considered "luxury" is not published.
Characteristics
Features
Luxury cars have traditionally emphasized higher levels of comfort and safety. Manufacturers often introduce new safety technologies and comfort amenities on luxury models before they are available on more affordable models. Some brands, like Audi and BMW have expanded their marketing by "introducing lesser priced and strip-down economy versions of their products."
Luxury vehicles can be a status symbol for conspicuous consumption. However, since many European luxury car buyers shy away from conspicuous consumption, brands offer buyers the option of removing exterior badges that identify the model name or engine size.
The suspension system of most luxury cars is tuned to prioritize ride quality over handling; however, some are marketed as "sports luxury" and have a greater emphasis on handling characteristics.
Layout and powertrain
Traditionally, luxury cars have used a front-engine, rear-wheel drive (FR) layout. The FR layout is more expensive to produce and produces lower fuel economy than a front-wheel drive layout; however, it allows for larger engines (particularly straight-six, V8, and V12) to be used.
Some American luxury cars during the 1970s through the 1990s switched to a front-wheel drive layout with a transverse engine provoked by the Arab Oil Embargo of 1973 and the 1979 oil crisis, which caused automakers to discontinue many FR platforms in favor of the more economical front-wheel drive (FF) layout. From the early 2000s, several of these American luxury cars reverted to FR layouts.
Since the introduction of the Bentley Continental GT in 2003, additional luxury grand tourers feature all-wheel drive.
History
European manufacturers
Prior to World War II, a wide array of European producers made luxury cars, including Rolls-Royce, Bugatti, Delage, Delahaye, Talbot-Lago, Bentley, Alvis, Avions Voisin, Isotta Fraschini, Horch, Simson, Stoewer, Maybach, Mercedes-Benz, Hispano Suiza, Daimler Company, and Spyker.
France was a leading producer of powerful luxury automobiles before World War II. After World War II, the French government used puissance fiscale tax regulations to encourage manufacturers to build cars with small engines and French motorists to buy them. The Maserati-powered Citroën SM and the Citroën C6 were arguably the last domestic French luxury cars. In the 2010s, some French manufacturers have attempted to develop luxury cars; however the lack of a historical legacy has hindered these efforts. In 2014, Citroën introduced DS Automobiles sub-brand to market luxury cars.
Pre World War II, intermediate car manufacturers like Renault, Fiat, Opel, Lancia, Škoda, Riley, Praga, Peugeot, Hillman, and Tatra made luxury cars. However, they had to transition to produce economy cars and superminis post World War II.
Following World War II, Germany rose to become an export powerhouse, building on success with the Mercedes-Benz brand. Aircraft engine manufacturer BMW, began making motorcycles, then small cars, including under license from the Austin Motor Company, and evolved into the luxury market segment ultimately acquiring Rolls-Royce Motor Cars in 1998. Volkswagen entered the high-end market by expanding or acquiring additional brands such as Audi, Porsche, Bentley, Lamborghini, and Bugatti.
In the Soviet Union, the manufacturer ZiL (then called Zis) began producing representational limousines in the mid-1930s. In the early 1950s, GAZ joined with the somewhat smaller "Chaika" model range. In 2018, Aurus Motors was established to produce luxury vehicles for the Russian market.
North American manufacturers
The luxury car phenomenon began at the start of the automobile industry when the wealthy frequently invested in manufacturing such models to gain social prestige. Emphasis was also placed on custom-built coachwork. The 1920s and 1930s were the apogee of production of these very large luxury automobiles from many manufacturers. The significant North American manufacturers from 1910 until 1940 included Auburn, Buick, Cadillac, Chrysler, Continental, Cord, Daniels, DeSoto, Duesenberg, Franklin, Imperial, LaFayette, LaSalle, Lincoln, Marmon, Packard, Peerless, Pierce Arrow, Ruxton, Stearns-Knight, and Stutz. The Great Depression put many luxury car manufacturers out of business; others would hold on before going defunct during the postwar era.
From 1946 until the late 1990s, Cadillac was the top-selling brand of luxury cars in the U.S., while Lincoln was second. The most successful and long-running model names during this era were the Cadillac DeVille, Lincoln Continental, and the Chrysler Imperial. The Lincoln Mark Series and Cadillac Eldorado were positioned in the personal luxury category, and competition between them continued into the 1990s.
The personal luxury car emerged into mass popularity and affordability as an America-specific category of popularly-priced cars made from the 1950s by the four domestic manufacturers (GM, Ford, Chrysler, and AMC) that reached peak popularity in the 1970s. The cars were stylized, mass-produced two-door coupés or convertibles, relying on standard components. These distinctively styled cars were targeting the needs of individual customers, not an entire family. The longest running model lines were the 1958-1997 Ford Thunderbird, 1956-1998 Lincoln Mark Series, and the 1967-2002 Cadillac Eldorado.
In 1990, American luxury brands dominated, with Cadillac selling over a quarter-million cars, and Lincoln had its best year ever at 231,660 units. However, the market was changing with an ever greater acceptance of smaller, more efficient imported luxury brands while at the same time, the domestic manufacturers were downsizing their models with product decisions that backfired on quality and brand respect.
Since the late 1990s, Japanese and German brands have sold the most luxury-type cars in the United States. However, the Cadillac Escalade has led the luxury SUV segment sales in the United States since its introduction in 1998, with the highest sales for 15 out of its first 20 years on the market.
In the 2000s, both Ford and General Motors produced luxury pickups: 2002-2013 Cadillac Escalade EXT, 2002-2003 Lincoln Blackwood, and 2006-2014 Lincoln Mark LT.
In the late 2000s, the Cadillac CTS and Cadillac DTS led to a resurgence in the brand's luxury sedans. The equivalent sedan from the Ford group, the 2008 Lincoln MKS, was also regarded as a significant improvement over previous models. In 2010, BMW was the best-selling luxury vehicle manufacturer by sales, with Audi and Mercedes-Benz the second and third highest selling luxury brands.
East Asian manufacturers
Chinese manufacturer Hongqi was launched in 1958, making it the oldest Chinese luxury car marque. Later newcomers joined taking advantage of the rise of electric powertrains, with NEV brands such as Nio in 2014, Lynk & Co in 2016, HiPhi in 2019, and Zeekr in 2021 producing luxury electric and hybrid vehicles.
Japanese manufacturers have been producing luxury cars since the 1950s, including the Toyota Crown (1955–present), Prince/Nissan Gloria (1959–2004), Nissan Cedric (1960–2015), Mitsubishi Debonair (1964–1998), Nissan President (1965–2010), Toyota Century (1967–present), Mazda Luce/929 (1969–1991), and Honda Legend (1985–2021).
Since the 1980s, overseas sales of Japanese luxury cars have increased, challenging traditional European luxury brands.
Several East Asian manufacturers have created sub-brands for the marketing of luxury cars. The first of these was the 1986 launch of Acura (a Honda sub-brand), followed by Lexus (Toyota) in 1989, Infiniti (Nissan) in 1989, and Genesis (Hyundai) in 2015.
2007–2008 financial crisis and the Great Recession
The time of the 2007–2008 financial crisis and the Great Recession was the first time since the Great Depression that the luxury car market suffered considerably, something not seen in previous economic downturns. Many such customers saw their net worth decline during this time. For example, some of the steepest drop-offs came at the high end, including the BMW 7 Series and Rolls-Royce Phantom, and in 2010 Mercedes-Benz dropped the price of the W212 E-Class. The unusually sharp decline in luxury car sales has led observers to believe that there is a fundamental shift and reshaping of the luxury automotive market, with one industry official suggesting that the marques no longer command the premiums that they used to and another saying that conspicuous consumption was no longer attractive in poor economic conditions. Additionally, mainstream brands have been able to offer amenities and devices such as leather, wood, and anti-lock brakes, previously found only on luxury cars, as the costs decline.
However, luxury vehicle sales remained relatively high compared to their non-luxury counterparts. This was aided by growing interest in luxury vehicles from emerging markets such as China and Russia.
Sales in the entry-level luxury segment remained strong throughout the GFC, due to prices being lowered to compete with well-equipped non-luxury cars. For example, in Canada, several luxury manufacturers set sales records in August 2009, due mostly to discounted pricing on entry-level luxury vehicles.
Brands
Some auto manufacturers market their luxury models using the same marque as the rest of their line. Others have created a separate marque (e.g. Lexus, launched by Toyota in 1989) or purchased one (e.g. Bentley, by Volkswagen in 1998).
Occasionally, a luxury car is initially sold under a mainstream marque and is later rebranded under a specific luxury marque (for example, the Hyundai's Genesis to Genesis G80 as well as the Citroën DS to DS 5).
For mass-produced luxury cars, sharing of platforms or components with other models is common, as per modern automotive industry practice.
Market categories
Compact executive / compact luxury
A compact executive car or a compact luxury car is a premium car larger than a premium compact and smaller than an executive car. In European classification, compact executive cars are part of the D-segment. In North American terms, close equivalents are "compact premium car", "compact luxury car", "entry-level luxury car" and "near-luxury car". Compact executive cars are usually based on the platform of a mid-size car (also known as large family car or D-segment), while some models may be based on a compact car (also known as small family car or C-segment).
Executive / mid-size luxury
Executive car is a British term for an automobile larger than a large family car. In official use, the term is adopted by Euro NCAP, a European organization founded to test for car safety. It is a passenger car classification defined by the European Commission.
Luxury saloon / full-size luxury sedan
The next category of luxury cars is known in Great Britain as a "luxury saloon" or "luxury limousine," and is known in the United States as a full-size luxury sedan, large luxury sedan, or flagship sedan. It is the equivalent of the European F-segment and the German Oberklasse segment.
Many of these luxury saloons are the flagship for the marque and include the newest automotive technology. Several models are available in long-wheelbase versions, which provide additional rear legroom and may have a higher level of standard features.
Examples of luxury saloons / full-size luxury sedans include the BMW 7 Series, Jaguar XJ, Cadillac CT6, Genesis G90, Audi A8, Mercedes-Benz S-Class, Lexus LS,Hongqi H9, Porsche Panamera and Maserati Quattroporte.
Ultra-luxury
Luxury cars costing over (as of 2007) can be considered as "ultra-luxury cars". Examples include the Rolls-Royce Phantom, Maybach 57 and 62, Hongqi L5, Bentley Mulsanne, Cadillac Celestiq, Toyota Century, and Aurus Senat. High-end sports cars which are targeted towards performance rather than luxury are not usually classified as ultra-luxury cars, even when their cost is greater than . The history of a brand and the exclusivity of a particular model can result in price premiums compared to luxury cars with similar features from less prestigious manufacturers. Ultra-luxury cars are usually selected as the official state car.
Grand tourer
Grand tourers are essentially high-performance luxury vehicles. These vehicles are generally two-door coupes and are made for long-distance driving, combined with the luxury of an executive car or full-size luxury car.
Luxury SUV / crossover
Long before the luxury SUV segment was defined and became popular, the 1966 Jeep Super Wagoneer was marketed at the time as a fully-equipped station wagon. It was the first off-road SUV to offer a V8 engine and automatic transmission along with luxury car trim and equipment. Standard features included bucket seating, a center console, air conditioning, a seven-position tilt steering wheel, a vinyl roof, and gold-colored trim panels on the body sides and tailgate. By the late 1970s, optional equipment included an electric sunroof. The 1978 Jeep Wagoneer Limited was the spiritual successor to the Super Wagoneer and was the first four-wheel drive car to use leather upholstery.
The Range Rover was released in 1970 as a two-door vehicle for off-road durability with few "creature comforts." A four-door version was added in 1981, and the model was pushed upmarket in 1983 by introducing an automatic transmission (Chrysler's A727 TorqueFlite) as an option. The Range Rover had a long-travel coil-spring suspension and an aluminium V8 engine.
In the mid-1990s, the SUV market expanded with new entrants. By the mid-1990s, the entry-level Ford Explorer and upscale Jeep Grand Cherokee were the market leaders for SUVs. The fastest-growing sector of this market was for the so-called luxury SUVs, which included the Jeep Grand Cherokee ... the Grand Cherokee's allure: "This vehicle is proof you can have a true off-road vehicle without giving up luxuries and amenities" with the Jeep providing a crucial new intangible factor for buyers—image.
The SUV models generated higher profit margins than passenger cars, and car manufacturers began introducing new luxury SUVs during the late 1990s. SUVs such as the 1995 Lexus LX, 1997 Mercedes-Benz M-Class, and 1998 Lincoln Navigator were the first SUVs produced by these luxury car brands. Some of these early luxury SUV models used unibody construction, becoming part of the trend moving away from the body-on-frame construction traditionally used by off-road vehicles.
During the mid-2000s, SUVs from luxury car brands grew by almost 40% in the United States to more than 430,000 vehicles (excluding SUV-only brands like Hummer and Land Rover), at a time when luxury car sales suffered a 1% decline, and non-luxury SUV sales were flat. By 2004, 30% of major luxury brands' U.S. sales were SUVs. Crossover SUVs became increasingly popular in the mid-2000s, and manufacturers also began to produce luxury versions of crossovers. The Lexus RX was the earliest luxury crossover on the market, and it has since been the best-selling luxury vehicle in the US. Some luxury crossovers are built on a platform shared with sedans or hatchbacks. For example, the Infiniti FX is based upon the same platform as the Infiniti G35 sedans and coupes. While early luxury crossovers released in the late 1990s have resembled traditional boxy SUVs, later crossovers, such as the Infiniti FX and BMW X6, have been designed with a sporting appearance.
Despite the increased popularity of crossover models, traditional luxury SUVs remain in production. Examples include the Lexus LX, Infiniti QX80, and Lincoln Navigator.
Research data from the mid-2000s suggested that luxury SUV buyers did not consider traditional luxury cars (e.g. sedans and coupes), therefore the SUV is becoming the key to bringing new customers into luxury dealerships.
Luxury car companies have increasingly introduced SUV or crossover models in the 2010s. For example, Rolls-Royce Cullinan, Bentley Bentayga, Aston Martin DBX, Maserati Levante, Lamborghini Urus, and Ferrari Purosangue. Some brands, such as Lincoln, have even moved to an all SUV and/or crossover lineup.
Luxury MPV
Manufacturers such as Mercedes-Benz, Toyota, Lexus, Buick, Hongqi, Zeekr and Volvo have marketed upscale luxury MPVs as luxury vehicles, mainly marketed for Asian markets.
Luxury MPVs generally have 3-rows of six or seven seats; however, range-topping flagship models may also offer a 2-rows option with four seats, which typically have more features than their cheaper counterparts. By the early 2020s, manufacturers have found additional strategies to improve technologies, such as new materials, new systems, and improving exteriors.
Examples of luxury MPV models include Mercedes-Benz V-Class, Lexus LM, Buick GL8, Hongqi HQ9, Toyota Alphard, Volvo EM90 and the Zeekr 009.
| Technology | Motorized road transport | null |
252110 | https://en.wikipedia.org/wiki/Snowy%20owl | Snowy owl | The snowy owl (Bubo scandiacus), also known as the polar owl, the white owl and the Arctic owl, is a large, white owl of the true owl family. Snowy owls are native to the Arctic regions of both North America and the Palearctic, breeding mostly on the tundra. It has a number of unique adaptations to its habitat and lifestyle, which are quite distinct from other extant owls. One of the largest species of owl, it is the only owl with mainly white plumage. Males tend to be a purer white overall while females tend to have more extensive flecks of dark brown. Juvenile male snowy owls have dark markings that may appear similar to females until maturity, at which point they typically turn whiter. The composition of brown markings about the wing, although not foolproof, is the most reliable technique to age and sex individual snowy owls.
Most owls sleep during the day and hunt at night, but the snowy owl is often active during the day, especially in the summertime. The snowy owl is both a specialized and generalist hunter. Its breeding efforts and global population are closely tied to the availability of tundra-dwelling lemmings, but in the non-breeding season, and occasionally during breeding, the snowy owl can adapt to almost any available prey – most often other small mammals and northerly water birds, as well as, opportunistically, carrion. Snowy owls typically nest on a small rise on the ground of the tundra. The snowy owl lays a very large clutch of eggs, often from about 5 to 11, with the laying and hatching of eggs considerably staggered. Despite the short Arctic summer, the development of the young takes a relatively long time and independence is sought in autumn.
The snowy owl is a nomadic bird, rarely breeding at the same locations or with the same mates on an annual basis and often not breeding at all if prey is unavailable. A largely migratory bird, snowy owls can wander almost anywhere close to the Arctic, sometimes unpredictably irrupting to the south in large numbers. Given the difficulty of surveying such an unpredictable bird, there was little in-depth knowledge historically about the snowy owl's status. However, recent data suggests the species is declining precipitously. Whereas the global population was once estimated at over 200,000 individuals, recent data suggests that there are probably fewer than 100,000 individuals globally and that the number of successful breeding pairs is 28,000 or even considerably less. While the causes are not well understood, numerous, complex environmental factors often correlated with global warming are probably at the forefront of the fragility of the snowy owl's existence.
Taxonomy
The snowy owl was one of the many bird species originally described by Carl Linnaeus in his landmark 1758 10th edition of Systema Naturae, where it was given the binomial name Strix scandiaca. The genus name Bubo is Latin for "horned owl" and scandiacus is Neo-Latin for "of Scandinavia". The former generic name Nyctea is derived from Greek meaning "night". Linnaeus originally described the different plumages of this owl as separate species, with the male specimens of snowy owls being considered Strix scandiaca and the likely females considered as Strix nyctea. Until recently, the snowy owl was regarded as the sole member of a distinct genus, as Nyctea scandiaca, but mtDNA cytochrome b sequence data shows that it is very closely related to the horned owls in the genus Bubo and the species is now thusly often considered inclusive with that genus. However, some authorities debate this classification, still preferring Nyctea. Often authorities are motivated to retain the separate genus on the grounds of osteological distinctions.
Genetic testing revealed a reasonably distinct genetic makeup for snowy owls, being about 8% genetically distinct from other Bubo owls, perhaps giving credence to those who count the species as separate under Nyctea. However, a fairly recent shared origin in evolutionary history has been illustrated through a combination of genetic study and fossil review and there is little, other than osteology of the tarsometatarsus, to outright distinguish the snowy owl from other modern species like the Eurasian eagle-owl (Bubo bubo). Genetic testing has indicated that the snowy owl may have diverged from related species at around 4 million years ago. Furthermore, it has determined that the living species genetically most closely related to the snowy owl is the great horned owl (Bubo virginianus). On a broader scale, owls in general have, through genetic materials, been determined to be a highly distinct group, with outwardly similar groups such as Caprimulgiformes revealed to not be at all closely related. Within the owl order, typical owls are highly divergent from barn-owls. Furthermore, the Bubo genus likely clustered at some point during the evolutionary process with other largish owls, such as Strix, Pulsatrix and Ciccaba, based on broad similarities in their voice, reproductive behaviors (i.e. hooting postures) and a similar number and structure of chromosomes and autosomes. A number, but not all, of extant typical owls seem to have evolved from an ancient shared common ancestor with the Bubo owls. In addition to the question of relationship of the traditional Bubo owls to the snowy owls, ongoing ambiguity of the relationship of other similarly large-sized owls has been persistent. These have sometimes either been included in the genus or within separate genera, i.e. the Ketupu or fish owls and the Scotopelia or fishing owls. Despite the adaptive distinctions, the grouping of these large owls (i.e. Bubo, snowy, fish and perhaps fishing owls) appears to be borne out via research of karyotypes.
The fossil history of snowy owls is fairly well documented despite some early confusion on how to distinguish the skeletal structure of the snowy owls from eagle-owls. It was determined that the snowy owl once was distributed much more widely and far farther to the south during the Quaternary glaciation when much of the Northern Hemisphere was in the midst of an ice age. Fossil records shows that snowy owls once could be found in Austria, Azerbaijan, Czechoslovakia, England, France, Germany, Hungary, Italy, Poland, Sardinia and Spain as well as in the Americas in Cape Prince of Wales, Little Kiska Island, St. Lawrence Island, and in Illinois. In the Late Pleistocene the range expanded southward even more so to Bulgaria (80,000–16,000 years, Kozarnika Cave, W Bulgaria). and much of the Italian Peninsula. Pleistocene era fossils from France, i.e. B. s. gallica, showed that the snowy owls of the time were somewhat bulkier (though still notably smaller than contemporary eagle-owls of the times, which were larger than the eagle-owls of today) and osteologically more sexually dimorphic in size than the modern form (9.9% dimorphism in favor of females in the fossils against 4.8% in the same features today). There are no subspecific or other geographical variations reported in the modern snowy owls, with individuals of vastly different origins breeding together readily due to their nomadic habits. Despite apparent variations in body size, environmental conditions are the likely variant rather than genetics. No evidence could be found of phylogeographic variation in snowy owls upon testing. Furthermore, the snowy owl appears to have a similar level of genetic diversity as other European owls.
Hybrids
Snowy owls are not known to interbreed with other owl species in the wild, and accordingly, no hybrids of snowy owls and other owl species have yet been sighted in the wild. However, a hobby falconer in Kollnburg, Germany, bred hybrids from a male snowy owl and a female Eurasian eagle-owl (Bubo bubo) in 2013. The two resulting male hybrid owls possessed the prominent ear-tufts (generally absent in snowy owls), general size, orange eyes, and the same pattern of black markings on their plumage from their Eurasian eagle-owl mother, while retaining the generally black-and-white plumage colours from their snowy owl father. The hybrids were dubbed "Schnuhus" from the German words for snowy owl and Eurasian eagle-owl (Schnee-Eule and Uhu, respectively). As of 2014, the hybrids had grown to maturity and were healthy.
Description
The snowy owl is mostly white. They are purer white than predatory mammals like polar bears (Ursus maritimus) and Arctic fox (Vulpes lagopus). Often when seen in the field, these owls can resemble a pale rock or a lump of snow on the ground. It usually appears to lack ear tufts but very short (and probably vestigial) tufts can be erected in some situations, perhaps most frequently by the female when she is sitting on the nest. The ear tufts measure about and consist of about 10 small feathers. The snowy owl has bright yellow eyes. The head is relatively small and, even for the relatively simply adapted hearing mechanism of a Bubo owl, the facial disc is shallow and the ear is uncomplicated. 1 male had ear slits of merely on left and on the right. Females are almost invariably more duskily patterned than like-age males. In mature males, the upper parts are plain white with usually a few dark spots on the miniature ear-tufts, about the head and the tips of some primaries and secondaries whilst the underside is often pure white. Despite their reputation for being purely white, only 3 out of 129 Russian museum specimens of adult males showed an almost complete absence of darker spots. The adult female is usually considerably more spotted and often slightly barred with dark brown on the crown and the underparts. Her flight and tail feathers are faintly barred brown while the underparts are white in base color with brown spotting and barring on the flanks and upper breast. In confusingly plumaged snowy owls, the sex can be determined by the shape of wing markings, which manifest as bars more so in females and spots in males. However, the very darkest males and the lightest females are nearly indistinguishable by plumage. On rare occasions, a female can appear almost pure white, as has been recorded in both the field and in captivity. There is some evidence that some of the species grow paler with age after maturity. One study's conclusions were that males were usually but not always lighter and that correctly aging is extremely difficult, sometimes individuals either get lighter, darker or do not change their appearance with age. On the other hand, with close study, it is possible to visually identify even individual snowy owls using the pattern of markings on the wing, which can be somewhat unique in each individual. After a fresh moult, some adult females that previously appeared relatively pale newly evidenced dark, heavy markings. On the contrary, some banded individuals over at least four years were observed to have been almost entirely unchanged in the extent of their markings. In another very pale owl, the barn owl (Tyto alba), the sexual dimorphism of spotting appears to be driven by genetics while, in snowy owls, environment may be the dictating factor instead.
The chicks are initially grayish white but quickly transition to dark gray-brown in the mesoptile plumage. This type of plumage camouflages effectively against the variously colored lichens that dot the tundra ground. This is gradually replaced by plumage showing dark barring on white. At the point of fledging, the plumage often becomes irregularly mottled or blotched with dark and is mostly solidly dark gray-brown above with white eyebrows and other areas of the face white. Recently fledged young can already be sexed to a semi-reliable degree by the dark marking patterns about their wings. The juvenile plumage resembles that of adult females but averages slightly darker on average. By their second moult fewer or more broken bars are usually evidenced on the wing. The extent of white and composition of wing patterns become more dimorphic by sex with each juvenile moult, culminating in the 4th or 5th pre-basic moult, wherein the owls are hard to distinguish from mature adults. Moults usually occur from July and September, non-breeding birds moulting later and more extensively, and are never extensive enough to render the owls flightless. Evidence indicates that snowy owls may attain adult plumage at 3 to 4 years of age, but fragmentary information suggests that some males are not fully mature and/or as fully white in plumage that they can attain until the 9th or 10th year. Generally speaking, moults of snowy owls occur more quickly than do those of Eurasian eagle-owls.
The toes of the snowy owl are extremely thickly feathered white, while the claws are black. The toe feathers are the longest known of any owl, averaging at , against the great horned owl which has the 2nd longest toe feathers at a mean of Occasionally, snowy owls may show a faint blackish edge to the eyes and have a dark gray cere, though this is often not visible from the feather coverage, and a black bill. Unlike many other whitish birds, the snowy owl does not possess black wingtips, which is theorized to minimize wear-and-tear on the wing feathers in the other whitish bird types. The conspicuously notched primaries of the snowy owl appear to give an advantage over similar owls in long-distance flight and more extensive flapping flight. The snowy owl does have some of the noise-canceling serrations and comb-like wing feathers that render the flight of most owls functionally silent, but they have fewer than most related Bubo owls. Therefore, in combination with its less soft feathers, the flight of a snowy owl can be somewhat audible at close range. The flight of snowy owls tends to be steady and direct and is reminiscent to some of the flight of a large, slow-flying falcon. Though capable of occasional gliding flight, there is no evidence that snowy owls will soar. It is said that the species seldom exceeds a flying height of around even during passage. While the feet are sometimes described as "enormous", the tarsus is in osteological terms relatively short at 68% the length of those of a Eurasian eagle-owl but the claws are nearly as large, at 89% of the size of those of the eagle-owl. Despite its relatively short length, the tarsus is of similar circumference as in other Bubo owls. Also compared to an eagle-owl, the snowy owl has a relatively short decurved rostrum, a proportionately greater length to the interorbital roof and a much longer sclerotic ring surrounding the eyes while the anterior opening are the greatest known in any owl. Owls have extremely large eyes which are nearly the same size in large species such as the snowy owl as those of humans. The snowy owl's eye, at about in diameter, is slightly smaller than those of great horned and Eurasian eagle-owls but is slightly larger than those of some other large owls. Snowy owls must be able to see from great distances and in highly variable conditions but probably possess less acute night vision than many other owls. Based on the study of dioptres in different owl species, the snowy owl was determined to have eyesight better suited to long-range perception than to close discrimination, while some related species such as great horned owls could probably more successful perceive closer objects. Despite their visual limits, snowy owls may have up to 1.5 times more visual acuity than humans. Like other owls, snowy owls can probably perceive all colors but cannot perceive ultraviolet visual pigments. Owls have the largest brains of any bird (increasing in sync with the size of the owl species), with the size of the brain and eye related less to intelligence than perhaps to increased nocturnality and predatory behavior.
Size
The snowy owl is a very large owl. They are the largest avian predator of the High Arctic and one of the largest owls in the world. Snowy owls are about the sixth or seventh heaviest living owl on average, around the fifth longest and perhaps the third longest winged. This species is the heaviest and longest winged owl (as well as the second longest) in North America, the second heaviest and longest winged owl in Europe (and third longest) but is outsized in bulk by about 3 to 4 other species in Asia. Despite being sometimes described as of similar size, the snowy owl is somewhat larger in all aspects of average size than the great horned owl while the similarly specialized taiga-dwelling great grey owl (Strix nebulosa), is longer in total length and of similar dimensions in standard measurements, but is shorter winged and much less heavy than the snowy owl. In Eurasia, the Eurasian eagle-owl is larger in all standards of measurements than the snowy owl not to mention two additional species each from Africa and Asia that are slightly to considerably heavier on average than the snowy owl. Like most birds of prey, the snowy owl shows reverse sexual dimorphism relative to most non-raptorial birds in that females are larger than males. Sexual dimorphism that favors the female may have some correlation with being able to more effectively withstand food shortages such as during brooding as well as the rigors associated with incubating and brooding. Females are sometimes described as “giant” whereas males appear relatively “neat and compact”. However, the sexual dimorphism is relatively less pronounced compared to some other Bubo species.
Male snowy owls have been known to measure from in total length, with an average from four large samples of and a maximum length, perhaps in need of verification, of reportedly . In wingspan, males may range from , with a mean of . In females, total length has been known to range from , with a mean of and an unverified maximum length of perhaps (if so they would have the second longest maximum length of any living owl, after only the great grey owl). Female wingspans have reportedly measured from , with a mean of . Despite one study claiming that snowy owl had the highest wing loading (i.e. grams per square cm of wing area) of any of 15 well-known owl species, more extensive sampling demonstratively illustrated that the wing loading of snowy owls is notably lower than Eurasian eagle- and great horned owls. The conspicuously long-winged profile of a flying snowy owl compared to these related species may cause some to compare their flight profile to a bulkier version of an enormous Buteo or a large falcon. Body mass in males can average from , with a median of and a full weight range of from six sources. Body mass in females can average from , with a median of and a full weight range of . Larger than the aforementioned body mass studies, a massive pooled dataset at six wintering sites in North America showed that 995 males averaged at while 1,189 females were found to average . Reported weights of down to for males and of for females are probably in reference to owls in a state of starvation. Such emaciated individuals are known to highly impaired and starvation deaths are probably not infrequent in winters with poor food accesses.
Standard measurements have been even more widely reported than length and wingspan. The wing chord of males can vary from , averaging from with a median of . The wing chord of females can vary from , averaging from with a median of . The tail length of males can vary on average from , with a full range of and a median of . The tail length of females can average from , with a full range of and a median of . Data indicates that slightly longer wing chord and tail lengths were reported on average in Russian data than in American research, however the weights were not significantly different in the two regions. Less widely taken measurements include the culmen, which can measure from with a median average of in males and in females, and the total bill length which is from , with an average in both sexes of . Tarsal length in males averages about , with a range of , and averages about , with a range of , in females.
Identification
The snowy owl is certainly one of the most unmistakable owls (or perhaps even animals) in the world. No other species attains the signature white stippled sparsely with black-brown color of these birds, a coloring which renders their bright yellow eyes all the more detectable, nor possesses their obvious extremely long feathering. The only other owl to breed in the High Arctic is the short-eared owl (Asio flammeus). Both species inhabit open country, overlap in range and are often seen by day, but the short-eared is much smaller and more tan or straw-colored in coloration, with streaked brown on chest. Even the palest short-eared owls conspicuously differ and are darker than the snowy owl; additionally the short-eared most often hunts in extended flights. More similar owls such as the Eurasian eagle-owl and the great horned owl attain a fairly pale, sometimes white-washed look in their northernmost races. These species do not normally breed nearly as far north as snowy owls but overlaps certainly do occur with snowy owls when the latter owl sometimes comes south in winter. However, even the most pale great horned and Eurasian eagle-owls are still considerably more heavily marked with darker base colors than snowy owls (the whitest eagle-owls are paler than the whitest great horned owls), possess much larger and more conspicuous ear tufts and lack the bicolored appearance of the darkest snowy owls. While the great horned owl has yellow eyes like the snowy owl, the Eurasian eagle-owl tends to have bright orange eyes. The open terrain habitats normally used by wintering snowy owls are also distinct from the typical edge and rocky habitats usually favored by the great horned and Eurasian eagle-owls, respectively.
Vocalizations
The snowy owl differ in their calls from other Bubo owls, with a much more barking quality to their version of a hooting song. Perhaps as many as 15 different calls by mature snowy owls have been documented. The main vocalization is a monotonous sequence that normally contains 2–6 (but occasionally more), rough notes similar to the rhythm of a barking dog: krooh krooh krooh krooh... The call may end with an emphatic aaoow, which is somewhat reminiscent of the deep alarm call of a great black-backed gull (Larus marinus). They will call mainly from a perch but also sometimes do so in flight. The krooh call of the male snowy owl may perform multiple functions such as competitive exclusion of other males and advertising to females. The calls of this species may carry exceptionally far in the thin air of Arctic, certainly over more than , and maybe even to as much away. The female has a similar call to male but can be higher-pitched and/or more guttural as well as single notes which are often disyllabic, khuso. Female snowy owls have also been known to utter chirps and high screaming notes, similar to those of the nestlings. Both sexes may at times give a series of clucking, squeals, grunts, hisses and cackles, perhaps such as in circumstances when they are excited. The alarm call is a loud, grating, hoarse keeea. Another raspier bark is recorded, sometimes called a "watchman's rattle" call, and may be transcribed as rick, rick, rick, ha, how, quack, quock or kre, kre, kre, kre, kre. A female attacking to protect her nest was recorded to let out a crowed ca-ca-oh call, whilst other owls attacking to protect the nest did a loud version of the typical call while circling before dropping down. They may also clap their beak in response to threats or annoyances. While called clapping, it is believed this sound may actually be a clicking of the tongue, not the beak. Though largely only vocal in the breeding season, leading to some erroneous older accounts describing the snowy owl as completely silent, some vocalizations have been recorded in winter in the northern United States. Initially, the young of the snowy owl have a high pitched and soft begging call which develops into a strong, wheezy scream at around 2 weeks. At the point when the young owls leave the nest around 3 weeks, the shrill squeals they emit may allow the mothers to locate them.
Distribution and habitat
Breeding range
The snowy owl is typically found in the northern circumpolar region, where it makes its summer home north of latitude 60° north though sometimes down to 55 degrees north. However, it is a particularly nomadic bird, and because population fluctuations in its prey species can force it to relocate, it has been known to breed at more southerly latitudes. Although the total breeding range includes a little over , only about have a high probability of breeding, i.e. breeding at no more than 3–9-year intervals. Snowy owls nest in the Arctic tundra of the northernmost stretches of Alaska, Northern Canada, and the Euro-Siberian region.
Between 1967 and 1975, snowy owls bred on the remote island of Fetlar in the Shetland Isles north of mainland Scotland, discovered by the Shetland RSPB warden, Bobby Tulloch. Females summered as recently as 1993, but their status in the British Isles is now that of a rare winter visitor to Shetland, the Outer Hebrides and the Cairngorms. Vagrant snowy owls have occasionally been found as far south as Lincolnshire. Older records show that the snowy owls may have once semi-regularly bred elsewhere in Shetland. They range in northern Greenland (mostly Peary Land) and, rarely in "isolated parts of the highlands", Iceland. Thence, they are found breeding at times across northern Eurasia such as in Spitsbergen and western and northern Scandinavia. In Norway, they normally breed in Finnmark and seldom down as far south as Hardangervidda and in Sweden perhaps down to the Scandinavian Mountains while breeding is very inconsistent in Finland.
They also range in much of northern Russia, including northern Siberia, Anadyr, Koryakland, Taymyr Peninsula, Yugorsky Peninsula, Sakha (especially the Chukochya River) and Sakhalin. Breeding has also been reported sporadically to the south in the Komi Republic and even the Kama River in southern Perm Krai. Although considered part of their regular range, the last breeding by snowy owls in the Kola Peninsula was in the early 1980s; similarly, breeding maps show the species in Arkhangelsk Oblast and the Pay-Khoy Ridge but no breeding records are known in at least 30 years in either. They range throughout most of the Arctic isles of Russia such as Novaya Zemlya, Severnaya Zemlya, New Siberian Islands, Wrangel Island, Commander and Hall Islands.
In North America, the breeding range has been known in modern times to include the Aleutians (i.e. Buldir and Attu) and much of northern Alaska, most frequently from the Arctic National Wildlife Refuge to Utqiaġvik, and more sporadically down along the coastal-western parts such as through Nome, Hooper Bay, the Yukon Delta National Wildlife Refuge, and rarely even south to the Shumagin Islands. The snowy owl may breed extensively in northern Canada, largely making its home in the Arctic Archipelago. Their Canadian breeding range can include broadly Ellesmere Island up to Cape Sheridan, north coastal Labrador, the northern Hudson Bay, perhaps all of Nunavut (especially the Kivalliq Region), northeastern Manitoba, most of northern mainland and insular Northwest Territories (including the delta of the Mackenzie River) and northern Yukon Territory (where breeding is mostly confined to Herschel Island). Since breeding and distribution is very small, local and inconsistent in northern Europe, northern Canada and northern Alaska represent the core part of the breeding range for snowy owls along with several parts of northern and northeastern/coastal Russia.
Regular wintering range
During the winter months, many snowy owls leave the dark Arctic to migrate to regions further south. Southern limits of the regular winter range are difficult to delineate given the inconsistency of appearances south of the Arctic. Furthermore, not infrequently, many snowy owls will overwinter somewhere in the Arctic through the winter, though seldom appear to do so in the same sites where they have bred. Due in no small part to the difficulty and hazardousness of observation for biologists during these harsh times, there is very limited data on overwintering snowy owls in the tundra, including how many occur, where they winter and what their ecology is at this season. The regular wintering range has at times been thought to include Iceland, Ireland and Scotland and across northern Eurasia such as southern Scandinavia, the Baltics, central Russia, southwestern Siberia, Sakhalin southern Kamchatka and, rarely, north China and sometimes the Altai Republic. In North America, they occasionally regularly winter in the Aleutian island chain and do so broadly and with a fair amount of consistency in much of southern Canada, from British Columbia to Labrador. Recent research has indicated that snowy owls regularly winter in several of the northern seas during wintertime, following the leads of sea ice as perching sites and presumably hunting mostly seabirds in polynyas. In February 1886, a snowy owl landed on the rigging of the Nova Scotia steamship Ulunda on the edge of the Grand Banks of Newfoundland, over from the nearest land. It was captured and later preserved at the Nova Scotia Museum. Surprisingly, some studies have determined that after a high lemming year in North America, a higher percentage of snowy owls were using marine environments rather than inland ones.
Irruptive range
Large winter irruptions at temperate latitudes are thought to be due to good breeding conditions resulting in more juvenile migrants. These result in irruptions occurring further south than the typical snowy owl range in some years. They have been reported, as well as in all northerly states in the contiguous states, as far south as Georgia, Kentucky, South Carolina, nearly all the Gulf Coast of the United States, Colorado, Nevada, Texas, Utah, California and even Hawaii. In January 2009, a snowy owl appeared in Spring Hill, Tennessee, the first reported sighting in the state since 1987. Also notable is the mass southern migration in the winter of 2011/2012, when thousands of snowy owls were spotted in various locations across the United States. This was then followed by an even larger mass southern migration in 2013/2014 with the first snowy owls seen in Florida for decades. The nature of irruptions is less well-documented in Eurasia, in part due to the paucity of this owl in the European side, but accidental occurrence, presumably during irruptions, has been described in the Mediterranean area, France, Crimea, the Caspian part of Iran, Kazakhstan, northern Pakistan, northwestern India, Korea and Japan. Stragglers may too turn up as far south as the Azores and Bermuda.
Habitat
Snowy owls are one of the best known inhabitants of the open Arctic tundra. Frequently, the earth in snowy owl breeding grounds is covered with mosses, lichens and some rocks. Often the species preferentially occurs in areas with some rising elevation such as hummocks, knolls, ridges, bluffs and rocky outcrops. Some of these rises in the tundra are created by glacial deposits. The ground is usually rather dry in tundra but in some areas of the southern tundra can also be quite marshy. Not infrequently, they will also use areas of varied coastal habitat, often tidal flats, as a breeding site. Breeding sites are usually at low elevations, usually less than above sea level, but when breeding to the south in inland mountains, such as in Norway, they may nest at as high as . Outside the breeding season, snowy owls may occupy nearly any open landscape. Typically wintering sites are rather windswept with meager cover. These open areas can include coastal dunes, other coastal spots, lakeshores, islands, moorlands, steppes, meadows, prairies, other extensive grasslands and rather shrubby areas of the Subarctic. These may be favored due to their vague similarity to the flat openness of the tundra. Manmade open sites are now perhaps even more used than natural ones, often agricultural fields and rangeland, as well as large areas of cleared forests. During irruption years when they are found in the northeastern United States, juveniles frequent developed areas including urban areas and golf courses, as well as the expected grasslands and agricultural areas that older birds primarily use. On the plains of Alberta, observed snowy owls spent 30% of their time in stubble-fields, 30% in summer fallow, 14% in hayfield and the remainder of the time in pasture, natural grasslands and sloughs. The agricultural areas, large untouched by the farmers in winter, may have had more concentrated prey than the others in Alberta. Perhaps the most consistently attractive habitat in North America to wintering snowy owls in modern times may be airports, which not only tend to have the flat, grassy characteristics of their preferred habitats but also by winter host a particular diversity of prey, both pests which rely on humans as well as wildlife attracted to the extensively grassy and marshy strips that dot the large airport vicinities. For example, Logan International Airport in Massachusetts has relatively one of the most reliable annual populations known in the United States in winter. All ages spend a fair amount of their time over water in the Bering Sea, the Atlantic Ocean and even the Great Lakes, mostly on ice floes. These marine and ocean-like freshwater areas were observed to account for 22–31% of habitat used in 34 radio-tagged American snowy owls over two irruptive years, with the tagged owls occurring a mean of from the nearest land (while 35–58% used the expected preferred habitats of grassland, pasture and other agricultural land).
Behavior
Snowy owls may be active to some extent at both day, from dawn to dusk, and night. Snowy owls have been seen to be active even during the very brief winter daytime in the northern winter. During the Arctic summer, snowy owls may tend to peak in activity during the twilight that is the darkest time available given the lack of full nightfall. Reportedly, the peak time of activity during summer is between 9:00 pm and 3:00 am in Norway. The peak time of activity for those owls that once nested on Fetlar was reported between 10:00 and 11:00 pm. According to one authority, the least active times are at noon and midnight. As days become longer near autumn in Utqiaġvik, the snowy owls in the tundra become more active around nightfall and can often be seen resting during the day, especially if it is raining. During winter in Alberta, snowy owls were tracked in the daytime, despite being also active at night (as they were deemed too difficult to track). In the study, they were most active during 8:00–10:00 am and 4:00–6:00 pm and often rested mostly from 10:00 am to 4:00 pm. The owls were perched for 98% of observed daylight and seemed to time their activity to peak times for rodents. The variation of activity is probably in correspondence with their primary prey, the lemmings, and like them, the snowy owl may be considered cathemeral. This species can withstand extremely cold temperatures, having been recorded in temperatures as low as minus 62.5 degrees Celsius with no obvious discomfort and also withstood a 5-hour exposure to minus 93 degrees Celsius but may have struggled with oxygen consumption by the end of this period. The snowy owl has perhaps the second lowest thermal conduction to the plumage on average of any bird after only the Adelie penguin (Pygoscelis adeliae) and rivals the best insulated mammals, such as Dall sheep (Ovis dalli) and Arctic fox, as the best insulated polar creature. Presumably as many as 7 rodents would need to be eaten daily to survive an extremely cold winter's day. Adults and young both have been seen to shelter behind rocks to shield themselves from particularly harsh winds or storms. Snowy owls often spending a majority of time on the ground, perched mostly on a slight rise of elevation. It has been interpreted from the morphology of their skeletal structure (i.e. their short, broad legs) that snowy owls are not well-suited to perching extensively in trees or rocks and prefer a flat surface to sit upon. However, they may perch more so in winter though do so only mainly when hunting, at times on hummocks, fenceposts, telegraph poles by roads, radio and transmission towers, haystacks, chimneys and the roofs of houses and large buildings. Rocks may be used as perches at times in all seasons. Though often relatively sluggish owls, like most related species, they are capable of sudden dashing movements in various contexts. Snowy owls can walk and run quite quickly, using outstretched wings for balance if necessary. This owl flies with fairly rowing wingbeats, occasionally interrupted by gliding on stretched wings. The flight is fairly buoyant for a Bubo owl. When displaying, the male may engage in an undulating flight with interspersed wingbeats and gliding in a slight dihedral, finally dropping rather vertically to the ground. They are capable for swimming but do not usually do so. Some seen to be swimming were previously injured but young have been seen to swim into water to escape predators if they cannot fly yet. They will also drink when unfrozen water is available. Snowy owl mothers have been observed to preen their young in the wild, while pairs in captivity have been observed to allopreen. In the period leading up to breeding, snowy owls switched regularly between searching (for nesting grounds) and loafing, often searching less when snow cover was less extensive.
Snowy owls will fight with conspecifics in all seasons occasionally but this is relatively infrequent during breeding and rarer still during winter. Dogfights and talon interlocking may ensue if the fight between two snowy owls continues to escalate. A study determined that snowy owls are able to orient the whitest parts of their plumage towards the sun, spending about 44% of time oriented as such during sunny days and much less on cloudy days. Some authors interpret this as a presumed signal to conspecifics, but thermoregulation could also be a factor. It is known that during winter in Alberta that female snowy owls are territorial towards one another and may not leave an area for up to 80 days but males are nomadic, usually only staying 1–2 days in an area (seldom to 3–17 days). The females spent on average seven times as long in a given area than did males. During threat displays, individuals will lower the front of the body, stretch the head low and forward, with partially extended wings and feathers on the head and raise their back. If continuously threatened or cornered, the posture in the threat display may become still more contoured and, if pressed, the owl will like back and attempt to slash with its large talons. The threat displays of males are generally more emphatic than those of females. Although snowy owls have been considered as semi-colonial, they do not appear to fit this mold well. Nesting sites can be loosely clustered but this is a coincidental response to concentrated prey and each pair tends to be somewhat intolerant of each other. During winter, snowy owls are usually solitary but some aggregations have been recorded, especially nearer the Arctic when more narrow food selection can lead to up to 20–30 owls gathering in an area of about . Congregations were also recorded in the winter in Montana, where 31–35 owls wintered in a area, owls mostly grouped in loose aggregations of 5–10 owls each or occasionally side-by-side or about apart. In extreme cases in Utqiaġvik, the owls may have exceptionally close active nests that may be down to only apart. Juvenile males appear to be especially prone to loose associations with one another, appearing to be non-territorial and able to hunt freely in front of one another. In a area in and around Utqiaġvik, productive years may have about 54 nests while none may be found in poor years. Utqiaġvik may have about 5 owls in early summer every , have a nest spacing of and the owls territory size is about . In Churchill, Manitoba, nest spacing averaged about . In Southampton Island in a year when the owls nested there, nest spacing averaged , with the closest two apart and density per nest was . In Nunavut, densities could go from 1 owl per in a productive year to 1 owl per in a poor year and from 36 nests in a area to none at all. Owl density on Wrangel Island in Russia was observed be a single bird each . The first known study of winter territories took place in Horicon Marsh where owls ranged from each. In Calgary, Alberta, mean territory size of juvenile females in winter was and adult females was . Wintering owls in central Saskatchewan were radio-monitored, determining that 11 males had an average range of , while that of 12 females was with the combined average being .
Migration
It is fair to say that the snowy owl is a partial, if fairly irregular, migrant, having a very broad but patchy wintering range. 1st year birds tend to disperse farther south in winter than older owls with males wintering usually somewhat more to the south than females of equivalent ages, adult females often wintering the farthest north. The snowy owl likely covers more ground than almost any other owl in movements but many complex individual variations are known in movements, and they often do not take the traditional north–south direction that might be assumed. Migratory movements appear to be somewhat more common in America than in Asia. A study of wintering owls in the Kola Peninsula determined that the mean date of arrival of owls was 10 November with a departure date of 13 April, covering an average of during the course of the wintering period and clustering where prey was more concentrated. Some variety of movements recorded each autumn and snowy owls winter annually in plains of Siberia and Mongolia and prairies and marshlands of Canada. The Great Plains area of southern Canada host wintering snowy owls about 2 to 10 times more frequently than other areas of the continent. Some weak correlation has made with individuals having some level of fealty to certain wintering sites. Wintering snowy owls, a total of 419, recorded in Duluth, Minnesota from 1974 to 2012 would occur in larger numbers in years where rats were more plentiful. The amount of individual returns among 43 Duluth-wintering owls was fairly low in subsequent winters (8 for 1 year, a small handful in the next few years, and 9 in non-consecutive years). Sometimes surveys appeared to reveal hundreds of wintering snowy owls on coastal sea ice during an irruptive year. Three siblings that hatched in same nest in Cambridge Bay were recovered in drastically different spots at least a year later: one in eastern Ontario, one in Hudson Bay and one in Sakhalin Island. A nestling banded in Hordaland was recovered to the northeast in Finnmark. In the Logan Airport, 17 of 452 owls were recorded to return, eleven the following year, three 2 years later, and then singles variously 6, 10 and 16 years later. A banded female from Utqiaġvik was recorded to migrate over along seacoast down to Russia, returning over and covering at least in total. Another banded young female from Utqiaġvik went to the same Russian areas, returned to Utqiaġvik and then onto Victoria Island, but did appear to breed, while another also covered a similar route but ended up nesting on Banks Island. Another female migrated to the Canada–United States border, then moved back to the Gulf of Alaska, then to winter in the same border areas and then finally to both Banks and Victoria Island. Snowy owls from the Canadian Arctic were monitored to have covered an average of in one autumn then covered an average of a year later. In late winter, owls from the same area were found to have covered a mean of of ground in the tundra and spent a mean of 108 days, apparently searching for a suitable nesting situation the entire time.
In no fewer than 24 winters between 1882 and 1988, large numbers have occurred in Canada and the United States. These were irruption years. Record breeding irruptive years were recorded in the winters of 2011–2012 and 2014–2015. In the 1940s, it was calculated that the mean gape in time between large irruptions was 3.9 years. Southbound movements as such are much more conspicuous after peak vole years, once thought to be separated by periods of around 3–7 years. However, more extensive research has weakened the argument that irruptions are entirely food-based and the data indicates that irruptive movements are far from predictable. This is because a statewide survey in Alaska found no statewide synchrony in lemming numbers. Therefore, rather than decline of lemmings, it is the successful productivity of several pairs that plays the role, resulting in a large number of young owls that then irrupt. However, the snowy owls cannot breed in high numbers unless lemmings are widely available on the tundra. This connection of irruptions to high years of productivity was confirmed in a study by Robillard et al. (2016). About 90% of the snowy owls seen in irruptive years from 1991 to 2016 that were ageable were identified as juveniles.
Diet and hunting
Hunting techniques
Snowy owls may hunt at nearly any time of the day or night, but may not attempt to do so during particularly severe weather. During the summer solstice, the owls appear to hunt during "theoretical nightfall". Night-vision devices have allowed biologists to observe that snowy owls hunt quite often during the extended nighttime during the northern winter. Prey are both taken and eaten on the ground. Like other carnivorous birds, Snowy owls often swallow their small prey whole. Strong stomach juices digest the flesh, while the indigestible bones, teeth, fur, and feathers are compacted into oval pellets that the bird regurgitates 18 to 24 hours after feeding. Regurgitation often takes place at regular perches, where dozens of pellets may be found. Biologists frequently examine these pellets to determine the quantity and types of prey the birds have eaten. When large prey are eaten in small pieces, pellets will not be produced. Larger prey is often torn apart, sometimes including removal of the head, with the large muscles, such as the humerus or breast, typically eaten first. The scattering of remains that results from the increment feeding on larger prey is thought to result in under-identification of them compared to smaller prey items. The aptitude for hunting by day, hunting from the ground and hunting in almost always completely open and treeless areas are the primary ways in which the snowy owl differs in hunting from other Bubo owls. Otherwise, the hunting habits are similar. It is thought, due to their less refined hearing compared to other owls, prey is usually perceived via vision and movement. Experiments indicate that snowy owls can detect prey from as far as away. Snowy owls generally use a rise or, occasionally, a perch while hunting. 88% of observed 34 hunts in Utqiaġvik were undertaken from an elevated watch-site (56% mounds or rises, 37% telephone poles). Their hunting style may recall that of buzzards, with the hunting owl sitting rather low and perching immobile for a long spell. Although their usual flight is a slow, deliberate downbeat on the broad, fingered wings, when prey is detected from their perch, flight may be undertaken with a sudden, surprisingly quick accelerated style with interspersed wing beats. In Utqiaġvik, snowy owls may most frequently engage in a brief pursuit hunting style. In high winds capable of keeping their bulk aloft, snowy owls may too engage in a brief hovering flight before dropping onto prey. When hunting fish, apparently, some snowy owls will hover in a style reminiscent of the osprey (Pandion haliaetus), although in at least one other case a snowy owl was observed to capture fish by lying on its belly upon a rock by a fishing hole. A dashing stoop or pounce down onto their prey, ending in a high-impact "wallop", is fairly commonly recorded. Another common technique is the "sweep", wherein they fly by and grasp the prey while continuing to fly. In winter, snowy owls have been shown to be able to "snow plunge" to capture prey in the subnivean zone, under at least of snow. Perhaps least frequently, snowy owls may pursue their prey on foot, in doing so never taking wing. Snowy owls have been known to capture night-migrating passerines and shorebirds, sometimes perhaps on the wing, as well as large and/or potentially dangerous birds that were caught in air by snowy owls during daylight. On the wing pursuits against other various other carnivorous birds are sometimes undertaken as well to kleptoparasitize the prey caught by the other birds. Few variations of hunting technique were observed in winter observations from Alberta, almost all of the hunts being with the sit-and-wait method (also known as still-hunts). Adult females in Alberta had a considerably better hunting rate than juvenile females. Much as in Alberta, in Syracuse, New York, 90% of 51 hunts were still-hunting, with the sweep variant used after perch departure in 31% of hunts and the pounce method in 45% of hunts. The Syracuse-wintering owls used tall perches, a mixture of manmade objects and trees of around high, in nearly 61% of hunts, while nearly 14% were from low perches (i.e. fence-posts, snow banks and scrap piles) about half as high as the tall perches and started from a ground position nearly 10% of the time. In Sweden, males hunted from a perch more so than did females and adults both focused on significantly smaller prey (small mammals) and may have had more success hunting than juvenile snowy owls. Some snowy owls can survive a fast for up to about 40 days off of fat reserves. These owls were found to have extremely thick subcutaneous fat deposits of and it is likely owls that overwinter in the Arctic rely heavily on these to survive during this scarce time, in combination with lethargic, energy-conserving behavior.
Snowy owls may not infrequently exploit prey inadvertently provided or compromised by human activities, including ducks injured by duck hunters, birds maimed by antenna wires, various animals caught in human traps and traplines as well as domestic or wild prey being bred or farmed by humans in enclosures. A wide variety of accrued reports show that the snowy owl that scavenging on carrion is not uncommon (despite having once been thought to be very rare in all owls), including instances of reindeer (Rangifer tarandus) body parts brought to nests and owls following polar bears to secondarily feed on their kills. Even huge marine mammals such as walrus (Odobenus rosmarus) and whales can be fed upon by these owls when the opportunity occurs. Snowy owls produce a pellet that in different areas averages a median of about , averaging up to in length as in Europe.
Prey spectrum
The snowy owl is primarily a hunter of mammals. Most especially, they often live off of the northerly lemmings. Sometimes other similar rodents like voles and mice can also be found frequently in the snowy owl's foods. It is R-selected, meaning that it is an opportunistic breeder capable of taking advantage of increases in prey numbers and diversity, despite its apparent specialization. Birds are commonly taken as well, and may regularly include passerines, northern seabirds, ptarmigan and ducks, among others. Sometimes infrequent consumption of other prey such as beetles, crustaceans and occasionally amphibians and fish is reported (of these only fish are known to have been identified to prey species). It also takes carrion outside the breeding season. All told, more than 200 prey species have been known to be taken by snowy owls around the world. Generally, like other large owls (including even bigger owls like the Eurasian eagle-owl), prey selection tends toward quite small prey, usually small mammals, but they can alternate freely with prey that is much larger than typical given the opportunity or even bigger than themselves, including relatively large mammals and several types of large bird of almost any age. One study estimated for the biomes of Alaska and Canada, mean prey sizes for snowy owls were , in western North America, the mean prey size was and in eastern North America was , while the mean prey size in northern Fennoscandia was similar (at ). The mean number of prey species for snowy owls per biome ranged from 12 to 28. The opportunistic nature of snowy owls has long been known during their primarily winter observed feeding habits (leading to their unpopular nature and frequent persecution well into the 20th century).
Summer diet
The snowy owl's biology is closely tied to the availability of lemmings. These herbivorous rodents are largish members of the vole clan that are the predominant mammal of the tundra ecosystem alongside the reindeer and probably make up the majority of the mammalian biomass of the ecosystem. Lemmings are key architects of the soil, microtopography and plant life of the entire tundra. In the American lower Arctic areas, brown lemming of the Lemmus genus are predominant and tend to be found in lower, wetter habitats (feeding by preference on grasses sedges and mosses) while collared lemmings of the Dicrostonyx genus were in more arid, often higher elevation habitats with heathland and ate by preference willow leaves and forbs. The southerly brown lemmings behave differently than more northern collared lemming type, increasing almost limitlessly within preferred habitat whereas the collared type tends to spread to suboptimal habitats and therefore does not appear reach the high regional densities of the brown. Authorities now generally agree that there appears to be no synchrony between the brown and collared lemmings and the feeding access of snowy owls is irregular as a result, but snowy owls can likely alternate between the two lemming types as one or the other increases as they nomadically use different parts of the Arctic. It is possible that the rare coincidental mutual peak of both lemming types within a year results in the erratic high productivity that results in irruptions. Within individual Arctic lemming species, historically, populations can vary in rough 4- to 5-year trends. As a result, in areas such as Banks Island, the breeding rate of snowy owls can vary within a decade by about tenfold. Weights of lemmings taken can range from on Baffin Island, while those taken in Utqiaġvik averaged in female and male lemming, respectively. It was estimated based on captive daily food intake that a snowy owl may consume about of lemmings a day, though other estimates using voles show a daily need for about . On Southampton Island, 97% of the diet was lemmings. A very similar number of lemmings (nearly 100%) were found over 25 years of study in Utqiaġvik, amongst 42,177 cumulative prey items. Of 76 lemmings that could be identified to sex at a cache, male lemmings were found in the cache twice as often as female lemmings. While initial findings indicated on Wrangel Island that female lemmings outnumbered males in prey remains, to the contrary osteology indicated that, like Utqiaġvik, males were more often taken. However, the slightly larger, slower-moving females may be preferred when available.
In some areas, snowy owls can breed where lemmings are uncommon to essentially absent. Even in Utqiaġvik, where the diet is quite homogenously based in lemmings, the hatching of passerines, shorebirds, seabirds and waterfowl can provide a key resource when lemmings are not found regularly and may be the only means by which the young can survive at such lean times. In the Nome, Alaska area, the locally nesting snowy owls reportedly switched from lemmings to ptarmigans when the latter's chicks hatched. A somewhat varying diet was also reported in Prince of Wales Island, Nunavut where 78.3% of the biomass was lemmings, with 17.8% from waterfowl, 3.3% from weasel and about 1% from other birds. In Fennoscandia, among 2,700 prey items only a third were Norway lemmings (Lemmus lemmus) and a majority were voles at 50.6%, probably largely the tundra vole (Microtus oeconomus). A more detailed glance at Finnish Lapland showed that amongst 2,062 prey items, 32.5% of the foods were Norway lemmings (though in some years the balance could range up to 58.1%), 28% were grey red-backed voles (Myodes rufocanus) and 12.6% were tundra voles, with birds constituting a very small amount of the prey balance (1.1%). In northern Sweden, a more homogenous diet was found with the Norway lemming constituting about 90% of the foods. In the Yamal Peninsula, 40% of the diet was collared lemmings, 34% were Siberian brown lemming (Lemmus sibiricus), 13% were Microtus voles and ptarmigan and ducks both constituting 8% and with other birds making up much of the remaining balance. In some parts of the tundra, snowy owls may opportunistically prey upon Arctic ground squirrels (Spermophilus parryii). In the Hooper Bay area (much farther south than they usually nest), various rodents, in highland areas, and waterfowl, in marshland, were taken while breeding. When historically breeding on Fetlar in Shetland, the main prey for snowy owls was European rabbits (Oryctolagus cuniculus), Eurasian oystercatcher (Haematopus ostralegus), parasitic jaegers (Stercorarius parasiticus) and Eurasian whimbrel (Numenius phaeopus), in roughly that order, followed by other bird species with most (rabbits and secondary birds) prey taken as adults but for the oystercatchers and jaegers which were taken largely as fully grown but only recently fledged juveniles. 22–26% of oystercatcher and jaeger young in the island were estimated to be taken by snowy owls.
Bird predation by nesting snowy owls is highly opportunistic. Willow (Lagopus lagopus) and rock ptarmigan (Lagopus muta) of any age are often fairly regular in the diet of breeding snowy owls but they cannot be said to particularly specialize on these. Evidence was found in the Yamal Peninsula that the snowy owls became the primary predator of willow ptarmigan and that the predation was so frequent, it may have been the cause of the change of their habitat usage to willow thickets by the local ptarmigan. The reliance on ptarmigan has caused some conservation trickle-down concern for the owls because ptarmigan are hunted in large numbers, with the hunters of Norway permitted to cull up to 30% of the regional population. In North America, avian prey on the breeding ground regularly varies from small passerines like snow buntings (Plectrophenax nivalis) and Lapland longspurs (Calcarius lapponicus) to large waterfowl like king (Somateria spectabilis) and common eider (Somateria mollissima) and usually the goslings but also occasionally adults of geese such as brants (Branta bernicla), snow geese (Anser caerulescens) and cackling geese (Branta hutchinsii). Drake eiders of often similar size to the owls themselves are not infrequently the largest prey amongst remains around the nest mound. One nest had the bodies of all eiders that attempting to nest in the vicinity around it. The threatened and declining Steller's eider (Polysticta stelleri) when nesting in the Utqiaġvik area would appear to avoid the vicinity of snowy owl nests when selecting their own nesting sites due to the predation risk. Intermediately sized seabirds are often focused on in lieu of available lemmings. Foods were studied intensively in Iceland. Among 257 prey items found with a total prey mass of , birds made up 95% of the diet. The leading prey were adult rock ptarmigan, at 29.6% by number and 55.4% by biomass and adult European golden plover (Pluvialis apricaria), at 10.5% by number and 7.2% biomass. The rest of the balance was largely other shorebirds, which were taken slightly more often as chicks than adults. Pink-footed geese (Anser fabalis) were taken in equal number as goslings and adults, with respectively estimated average weights at these ages of . On the isle of Agattu, the diet consisted entirely of birds, as there are no mammals found there. The much favored food in Agattu was the ancient murrelet (Synthliboramphus antiquus), at 68.4% of the biomass and 46% by number, while the secondary prey were followed numerically by smaller Leach's storm-petrels (Oceanodroma leucorhoa) (20.8%) and Lapland longspurs (10%) and in biomass by smallish ducks, the green-winged teal (Anas carolinensis) and harlequin duck (Histrionicus histrionicus) (13.4% biomass collectively). In the Murman Coast of Russia, also in the absence of lemmings, seabirds formed the largest part of the diet.
Winter diet
On the wintering grounds, mammals often predominate in the snowy owl's food inland doing so less in coastal areas. Overall wintering snowy owls eat more diverse foods they do whilst breeding, furthermore coastal wintering snowy owls had more diverse diets than inland ones. As in summer, moderately sized water birds such as teal, northern pintail (Anas acuta) and numerous alcids and the like are often focused on when hunting birds. The diet in 62 pellets, amongst at least 75 prey items, from coastal Oregon showed the main foods as black rat (Rattus rattus) (at an estimated 40%), red phalarope (Phalaropus fulicarius) (31%) and bufflehead (Bucephala albeola) (19%). Witnessed attacks were mostly upon buffleheads in Oregon. In coastal southwestern British Columbia, the diet among 139 prey items was 100% avian. The predominant prey were water birds, mostly snatched directly from surface of the water and largely weighing , i.e. buffleheads (at 24% by number and 17.4% by biomass of foods) and horned grebes (Podiceps auritus) (at 34.9% by number and 24.6% by biomass), followed by variously other water birds, often the slightly larger species of glaucous-winged gull (Larus glaucescens) and the American wigeon (Mareca americana). A different study of this area also showed the predominance of ducks and other water birds to wintering snowy owls here, although Townsend's vole (Microtus townsendii ) (10.65%) and snowshoe hare (Lepus americanus) (5.7%) were also notably in a sample of 122 prey items.
During winter, snowy owls consume more strongly nocturnal prey than lemmings such as Peromyscus mice and northern pocket gophers (Thomomys talpoides). In southern Alberta, 248 prey items were found with North American deermouse (Peromyscus maniculatus), at 54.8% by number, and meadow voles (Microtus pennsylvanica), at 27% by number, as the main foods of snowy owls over two years. Other prey in Alberta were grey partridge (Perdix perdix) (at 5.79% of total), jackrabbits, weasels and owls. Richardson's ground squirrels (Urocitellus richardsonii) were consumed heavily in the Alberta study in a brief converged times of hibernation emergence and overwintering snowy owls. The sexual dimorphism in prey selection was also studied here, with male owls mainly focusing exclusively on the small rodents, females also took the same rodents but supplemented the diet with all alternate and larger prey. Overall, the meadow and montane voles (Microtus montanus) constituted 99% of over 4500 prey items in Montana. In Horicon Marsh in winter, 78% of the diet was meadow vole, with 14% being muskrats (Ondatra zibethicus), 6% ducks and smaller balances of rats and other birds. Snowy owls found in Michigan took meadow voles for 86% of the diet, white-footed mouse (Peromyscus leucopus) for 10.3% and northern short-tailed shrew (Blarina brevicauda) for 3.2%. Of 127 stomachs in New England in four irruptive winters from 1927 to 1942, of 155 prey items, 24.5% were brown rats, 11.6% were meadow voles and 10.3% were dovekie (Alle alle), with a smaller balance of snowshoe hare and birds from snow buntings to American black ducks (Anas rubripes). During the same years, stomach contents in Ontario included 40 identified prey items, led by brown rats (20%), white-footed mice (17.5%) and meadow voles (15%); of 81 prey items from Pennsylvania in 60 stomachs that were not empty, eastern cottontail (Sylvilagus floridanus) (32%), meadow vole (11.1%), domestic chicken (Gallus gallus domesticus) (11.1%) and northern bobwhite (Colinus virginianus) (5%) were the most often identified prey species. Introduced common pheasants were found to be somewhat more vulnerable than native American gamebirds like ruffed grouse due to their tendency to crouch rather than flush when approached by a flighted predator like the snowy owl in a glade or field. Some snowy owls wintering on rocky coasts and jetties were known in New England to live almost entirely off of purple sandpipers (Calidris maritima). The availability of brown rats may draw snowy owls to seemingly unattractive settings such as garbage dumps and under bridges. Meanwhile, snowy owls wintering in Lowell, Massachusetts were seen to live largely off of rock doves (Columba livia) caught off of buildings. Of 87 prey from stomachs in Maine, 35% were rats or mice, 20% were snowshoe hares and 10% were passerines. A small study of 20 prey items in an irruptive winter in Kansas found that 35% of the prey were red-winged blackbird (Agelaius phoeniceus), 15% prairie voles (Microtus ochrogaster) and 10% each by American coot (Fulica americana) and hispid cotton rats (Sigmodon hispidus).
On the isle of St. Kilda, 24 pellets were found for non-breeding snowy owls that stayed through the early summer. Of 46 prey items, the St Kilda field mouse (Apodemus sylvaticus hirtensis) was predominant by number at 69.6% but constituted 16.8% of biomass while adult Atlantic puffin (Fratercula arctica) constituted 63.5% of the prey biomass and 26% by number (rest of the balance being juvenile puffins and great skuas (Stercorarius skua)). The main subspecies of wood mouse was similarly dominant in the diet within County Mayo, Ireland and were presumably snatched at night due to their strict nocturnality. In Knockando, the winter diet was led by European rabbits (40.1%), red grouse (Lagopus lagopus scotica) (26.4%) and adult mountain hare (Lepus timidus) (20.9%) (in 156 pellets); in Ben Macdui, the diet was led by rock ptarmigan (72.3%), field voles (Microtus agrestis) and juvenile mountain hare (8.5%) (33 pellets); in Cabrach, the diet was led by red grouse (40%), mountain hare (20%) and European rabbit (15%) (16 pellets). Among 110 prey items found for snowy owls found wintering during irruption in southern Finland, all but 1 prey item were field voles (the only other prey being a single long-tailed duck (Clangula hyemalis)). Far to the east, wintering owls in the Irkutsky District were found to subsist mostly on narrow-headed voles (Microtus gregalis). In a wintering population in Kurgaldga Nature Reserve of Kazakhstan, the main foods were grey red-backed voles at 47.4%, winter white dwarf hamster (Phodopus sungorus) at 18.4%, steppe pika (Ochotona pusilla) at 7.9%, muskrat at 7.9%, Eurasian skylark (Alauda arvensis) at 7.9%, grey partridge at 5.3%, and both steppe polecat (Mustela eversmanii) and yellowhammer (Emberiza citrinella) at 2.6%. On the Kuril Islands, wintering snowy owls' main foods were reported as tundra voles, brown rats, ermines, and whimbrel, in roughly that order.
While most of the prey species are relatively small, snowy owls can prey on a fairly diverse size of birds and mammals. Data from the Logan Airport in over 6,000 pellets shows that meadow vole and brown rats predominated the diet in the area, supplanted by assorted birds both small and large. American black ducks were primarily taken among bird species with other birds taken here including relatively large and diverse species up to the size of Canada geese (Branta canadensis) and great blue herons (Ardea herodias) along with brants, American herring gulls (Larus argentatus), double-crested cormorants (Phalacrocorax auritus). Additionally, other birds as large as western capercaillie (Tetrao urogallus) (of both sexes), greater sage-grouse (Centrocercus urophasianus), yellow-billed loons (Gavia adamsii) and cygnets of Bewick's swans (Cygnus columbianus bewickii) can be taken by snowy owl. Among large mammalian prey species, snowy owls prey on both young and adults of large leporids such as Arctic hare (Lepus arcticus), Alaskan hare (Lepus othus), mountain hare and white-tailed jackrabbits (Lepus townsendii). At the other end of the scale, the snowy owl has been known to take birds down to the size of dark-eyed juncos (Junco hyemalis) and mammals down the size of common shrews (Sorex araneus). Fish are rarely taken anywhere but the snowy owl has been known to prey upon Arctic char (Salvelinus alpinus) and lake trout (Salvelinus namaycush).
Interspecific predatory relationships
The snowy owl is in many ways a very unique owl and differs from other species of owl in its ecological niche. Only one other owl, the short-eared owl, is known to breed in the High Arctic. However, the snowy owl shares its primary prey, the brown and collared lemmings, with a number of other avian predators. In sometimes differing parts of the Arctic, competing predators for lemmings are, in addition to short-eared owls, pomarine jaegers (Stercorarius pomarinus), long-tailed jaegers (Stercorarius longicaudus), rough-legged buzzards (Buteo lagopus), hen harriers (Circus cyaenus), northern harriers (Circus hudsonius) and generally less specialized gyrfalcons (Falco rusticollis), peregrine falcons (Falco peregrinus), glaucous gulls (Larus hypoboreus) and common ravens (Corvus corax). Certain carnivorous mammals, especially the Arctic fox and, in this region, the ermine, are also specialized to hunt lemmings. Most of the lemming predators are intolerant of the competition given the scattered nature of lemming populations and will displace and/or kill one another given the chance. However, given the need to conserve energy in the extreme environment, the predators may react passively to one another. When unusually breeding south in the Subarctic such as western Alaska, Scandinavia and central Russia, the number of predators with which the snowy owls are obligated to share prey and compete with may be too numerous to name. The taking of the young and eggs of snowy owls has been committed by a large number of predators: hawks and eagles, the northern jaegers, peregrine and gyrfalcons, glaucous gulls, common ravens, Arctic wolves (Canis lupus arctos), polar bears, brown bears (Ursus arctos), wolverines (Gulo gulo) and perhaps especially the Arctic fox. Adult snowy owls on the breeding grounds are far less vulnerable and can be justifiably qualified as an apex predator. Instances of the killing of adult snowy owls on the breeding grounds have been witnessed to be committed by a pair of pomarine jaegers on an incubating adult female snowy owl (possibly merely a competitive attack as she was left uneaten) and by an Arctic fox that killed an adult male snowy owl.
When it goes south to winter outside of the Arctic, the snowy owl has the potential to interact with a number of additional predators. By necessity, it shares its wintertime diverse prey with a number of formidable predators. These are known to include their cousins, the great horned owl and the Eurasian eagle-owl. They are relieved of heavy competition from the related species by differing temporal activity, i.e. being more likely to actively hunt in the daytime, and by habitat, using rather more open (quite often nearly treeless) habitats than them. During a study of wintering snowy owls in Saskatchewan, the authors indicated that the snowy owls may avoid areas inhabited and defended by great horned owls. Although they usually occurred here outside of an radius of central great horned owl ranges, they did not avoid the radius and different habitat usage may be a dictating factor. Given their mildly slighter size, it is unlikely that great horned owls (unlike the larger eagle-owl) would regularly dominate snowy owls in interactions and either species may give way to others depending on the size and disposition of the owls involved. Little study has been undertaken into the trophic competition of snowy owls with other predators during winter and, due to their scarcity, few predators are likely to expel much energy on competitive interactions with them, although many other predators will engage in anti-predator mobbing of snowy owls. Largely in winter, snowy owls have been the victim of a number of larger avian predators, though attacks are likely to be singular and rare. Instances of predation on snowy owls are known to have been committed several times in winter only by Eurasian eagle-owls. Additionally, golden eagles (Aquila chrysaetos) have been known to prey on snowy owls as well as all northern sea eagles: the bald (Haliaeetus leucocephalus), white-tailed (Haliaeetus albicilla) and Steller's sea eagles (Haliaeetus pelagicus). Snowy owls are also sometimes killed by birds that are mobbing them. In one instance, a peregrine falcon killed a snowy owl in a stoop after the owl had killed a fledgling falcon. Anecdotal reports indicate predation by gyrfalcons (on snowy owls of unknown age and condition) but it was possibly also an act of mobbing. In another, a huge throng of Arctic terns (Sterna paradisaea) relentlessly swarmed and attacked a snowy owl until it met its demise.
Almost certainly more often than being a victim of other predators, snowy owls are known to dominate, kill and feed on a large diversity of other predators. Snowy owls, much like other Bubo owls, will opportunistically kill other birds of prey and predators. Although they will readily plunder the nests of other raptorial birds given the opportunity, most predations are on full-grown raptorial birds during winter due to the scarcity of raptor nests in the open tundra. In addition, most competing predators of the Arctic, except the very large mammals, are probably vulnerable to a hungry snowy owl. In data from the Logan Airport alone over different winters, the snowy owls were observed to have preyed upon an impressive diversity of other raptorial birds: rough-legged buzzards, American kestrels (Falco sparverius), peregrine falcons, barn owls, other snowy owls, barred owls (Strix varia), northern saw-whet owls (Aegolius acadicus) and short-eared owls. While owls are likely encountered during corresponding hunting times, it is likely that the swift falcons are usually ambushed at night (much as other Bubo owls will do). In both the tundra and the wintering ground, there are several accounts of predation by snowy owls on short-eared owls. In addition, snowy owls have been known to prey on northern harriers, northern goshawks (Accipiter gentilis) and gyrfalcons. Snowy owls are also capable of taking large mammalian carnivores. Both juvenile and adult Arctic foxes have been known to fall prey to snowy owls, while predation by snowy owls on red foxes was reported in the Irkutsky District of Russia. A wintering snowy owl in Saskatchewan was observed to have preyed on an adult red fox (Vulpes vulpes) weighing around which may be the largest known prey known for snowy owls. Other relatively large carnivoran prey include adult house cat (felis catus), American mink (Mustela vision), and striped skunk (Mephitis mephitis). Also, several members of the weasel family, both small and relatively large, are known to be opportunistically hunted by snowy owls. As a result of its potential predator status, the snowy owl is frequently mobbed at all times of the year by other predatory birds, including fierce dive-bombing by several of the northern falcons on the wintering grounds, including even by the relatively tiny but fierce and very agile merlin (Falco columbarius). The much bulkier snowy owls cannot match the speed and flight ability of a falcon and may be almost relentlessly tormented by some birds such as peregrines.
Breeding
Pair bond and breeding territory
In Utqiaġvik, of 239 recorded breeding attempts, 232 were monogamous, the other 7 social bigamy. On Baffin Island, 1 male bred with 2 females and sired 11 total fledged young. Another case of bigamy was reported in Norway where the 2 females bred to one male were apart in nest site location. On Feltar from 1967 to 1975, a male bred with two females, 1 younger and was possibly his own daughter. In the Feltar males first time breeding with both females, he did not bring food to the younger female. However, when older female disappeared the following year, the male and younger female producing 4 young, but disappeared the subsequent year altogether in 1975. There are also unconfirmed cases of polyandry, with 1 female being fed by 2 males. Snowy owls can breed once per year but when food is scarce many do not even attempt to breed. Despite frequent wandering in search of food, they generally adhere more to a strict breeding season than short-eared owls nesting in the tundra. 9 radio-tagged female snowy owls about Bylot Island were tracked to study how pre-laying snow cover affects their searching behavior for breeding areas. These tracked females searched an average of 36 days and covered an average of . It is thought that the male and female mutually find an attractive breeding spot independently and converge. The breeding territory normally averages about as in both Baffin Island and Ellesmere Island but varies in accordance to abundance of food and density of owls. Nesting territories average at Baffin island in the range of during poor lemming years. Nesting territories may up to on Southampton Island and had a mean distance of between active nests. In Utqiaġvik, nesting pairs can vary from none to at least 7 and the territories average , with mean nest distances of . In the Norwegian highlands, nesting occurs only at times of plenty distances of between nests, averaging . Males marks territory with singing and display flights and likely always initiates. During the display, he engages in exaggerated wing beats with a shallow undulating and bouncy courtship flight with wings held in a dihedral. He often drops to the ground but then flies again to only glide gently back down. Overall, the flight is somewhat reminiscent of the flight of a moth. Females will answer her mate with her song during courtship. While courting, the male often also carries a lemming in his bill, then bows with cocked tail, similarly as in related owls (seldom displaying some other prey like snow buntings). He then flaps his wings open in an emphatic manner, with the ground display being relatively brief (about 5 minutes). The female may possibly refuse to breed if ritual not performed. A possible courtship was engaged in by a male in southern Saskatchewan when a female was sighted. On Southampton Island, at least 20 males observed in late May in a "lemming year". Nesting territory defense displays, not highly different from courtship displays, includes undulating flight and stiffly raised wings with bouts of exaggerated, delayed wing beats, looking like enormous white moths exposing their white wings under the sun. At times, competing males will interlock claws in mid-air. Territorial and nuptial displays are followed by a ground display by the male with the wings arched up in an "angel" posture, visible for well over a mile.
Nest sites
Most individuals arrive at the nest site by April or May with a few overwintering arctic exceptions. Males advertises potential nest sites to his mate by scratching the ground and spreading his wings over it. The nest is usually a shallow depression on a windswept eminence in the open tundra. There seems to be a variety of qualifiers for appropriate nest sites. The nest site is typically snow-free and dry relative to the surrounding environment, usually with a good view of the surrounding landscape. The nest may be made of ridges, elevated mounds, high polygons, hummocks, hills, human-made mounds and occasionally rocky outcrops. If covered with vegetation, taller plants that may obstruct view are plucked away sometimes. The nest sites are often long-established and naturally created by the freeze-thaw process of the tundra. Gravel bars may be used as well. The female may take the most active role in the nest's condition of any owl species. No owl build their own nests but female snowy owls take about three days constructing a scrape, digging with her claws and rotating until a fairly circular bowl is formed. She will still not construct or add foreign materials to the nest (despite some circumstantial evidence of moss and grass from outside the nest mound being found). In two separate cases in Utqiaġvik, two separate females dug out a second scrape to the side and below the main nests and appeared to have called all chicks to the more secluded nest to ride out severe weather until the skies cleared. The Utqiaġvik nest scrapes averaged in 91 with a mean depth of while the scrapes were smaller in Hooper Bay, reportedly diameter and in depth. Occasionally, in the lower tundra, snowy owls may too use old nests of rough-legged buzzards as well as abandoned eagle nests. Unlike other northerly breeding raptorial birds, the snowy owl is not known to nest on cliffs and the like, so do not enter into direct competition with eagles, falcons, ravens or other Bubo owls when nesting to the relative south. The area of nest mound often has a relatively rich plant life which attract the lemmings, which may tunnel right under and around the owl's nest. Geese, ducks and shorebirds of several species known to gain incidental protection by nesting close to snowy owls. Conversely, the snowy owls will sometimes kill and eat both young and adults of these birds, which implies a trade-off in the benefits.
Eggs
Egg-laying normally begins during early May to the first 10 days of June. Late thaws are harmful to them since they allow too little time for the full breeding process, with particularly importance given to good food supply in May for adults, even more so apparently than food supply in July when young are being fed. Late nests are possible cases of inexperienced pairs, low food supplies, bigamy or even replacement clutches. The clutch is extremely variable in size averaging around 7–9, with up to 15 or 16 eggs recorded in extreme cases. The clutch size very large relative to related species. Mean clutch sizes were 7.5 in a sample of 24 in Hooper Bay (range of 5–11); 6.7 in a sample of seven from Utqiaġvik (4–9); 9 in a sample of a sample of 5 in Baffin Island; 9.8 on Victoria Island; 8.4 (in a sample of 14) on Elsemere Island; 7.4 on Wrangel Island and 7.74 in Finnish Lapland. The average clutch size was 9.8 in a good year in Victoria Island while in a good year in Utqiaġvik the mean was 6.5. The clutch is laid directly to the ground and are pure, glossy white. An average egg is around with a range of heights from and diameter of . Egg weights are around , the median or average being in different datasets. The average egg size is relatively small, about 20% smaller than Eurasian eagle-owl eggs and 8% smaller than great horned owl eggs. Laying intervals are normally 2 days (41–50 hours mostly). The laying intervals can range up to 3–5 days in inclement weather. The laying of a clutch of 11 eggs can take 20–30 days, while a more typical nest of around 8 takes about up to 16 days. The interval between the 8th and 9th eggs can be up to about 4 days. Incubation begins with the first egg and is by female alone, while she is fed by her mate.
Parental behavior
Food is brought to the nest by males and surplus food is stored nearby. Females in breeding season often develop a very extensive brood patch which in this species is a fairly enormous, high vascularized featherless area of pink belly skin. Incubation lasts 31.8–33 days (unconfirmed and possibly dubious reports from as little as 27 to as much as 38-day incubations). The female alone broods the young, often while simultaneously incubating still unhatched eggs. Sometimes older chicks incidentally brood their younger siblings and females may shelter the young under her wings during inclement weather. When first feeding the young, the female may dismantle prey to feed the young only the softer body parts then gradually ramping up the size of proportions until they eat a whole prey item. Aggressive encounters with parent snowy owls are said to be "genuinely dangerous" and one resource claimed the snowy owl to be the bird species with the most formidable nest defense displays towards humans. The usual response to sighted humans near the nest is mild but continued approach begins to increasingly irritate the parents. At times, humans are forcefully dive-bombed upon, while other potential threats are dealt with in a “forward-threat” where the male walks towards the intruders, engaging in impressive feather-raising and fanning out of half-spread wings until they run forward and slash with both their feet and bill. Fairly serious injuries have been sustained in the worst of snowy owl defensive attacks, including cranial trauma, requiring researchers to make the long trek back to medical care, although human fatalities are not known. Snowy owl parents have been seen to aggressively attacked glaucous gulls, arctic fox and dogs in breeding ground in Utqiaġvik. Non-predatory animals like caribou in Utqiaġvik and sheep (Ovis aries) in Fetlar are attacked as well, possibly to avoid potential trampling of the eggs or the young. Males are said to do the majority of nest defense but the female will also often become involved as well. Analysis showed in Lapland, Sweden, that females in nest defense against people engaged in vocal displays (warning and mewing calls) and that males did not engage in mewing but did engage in most hooting calls, many warning calls and almost all physical attacks. In other instances, distraction displays are engaged in against predators, with a "broken-wing act" including high, thin squeals interspersed with weird squeaks, often taking flight only to quickly fall from the sky and imitate a struggle. One author recorded a male to draw him about from the nest before ceasing. 77% of 45 distraction displays in Lapland, Sweden were by females.
Development of young
Hatching intervals are generally from 1 to 3 days, quite often within 37–45 hours apart. New chicks are semi-altricial (i.e. typically helpless and blind), initially being white and rather wet but dry by the end of the first day. The weight of 7 hatchlings was , with an average of while 3 were . Due to the pronounced asynchrony of the egg-laying and hatching, the size difference between siblings can be enormous and in some cases when the smallest chick weighs only , the biggest chick already has attained a weight of around . When the oldest chick is about 3 weeks, the female will start to hunt as well as the male and both may directly feed the young although in some cases they may not need hunt very much if lemmings are particularly numerous. Caches of lemmings around a nest may include more than 80 lemmings that can support the family. Unlike many owls, the chicks of snowy owls are not known to behave aggressively toward one another or to engage in siblicide, perhaps in part due to the need for energy conservance. Some cases of cannibalism of chicks by the family group were thought to be cases where chicks die from other causes. When they are about 2 weeks, the chicks may begin to walk around the nest site which they leave by 18–28 days, although they are still unable to fly and may find safety in nooks and crannies of vegetation and rocks usually only about from the nest mound, as well as via their parents defense. Leaving the nest is thought to likely be an anti-predator strategy. The male snowy owl may drop fresh prey deliveries directly on the ground near the wandering young. After about three weeks of age, the young may wander fairly widely, rarely to , but usually stay within of the nest mound. Threat postures by young in reaction to researchers were first noticeable at about 20–25 days of age and common at about 28 days and the chicks can be impressively quick and agile-footed. The first fledgling occurs at around 35–50 days, and by 50–60 days the young can fly well and hunt on their own. The total care period is for 2–3.5 months, increasing in length with increased size of the brood. Although independence was once thought to be sought by late August or early September but is more likely by late September to October when migration season for the species begins. The nesting cycle is similar in length to the Arctic short-eared owls and faster than Eurasian eagle-owls by up to 2 months.
Maturity and nesting success
Sexual maturity reached the following year but the first breeding is normally at no sooner than the end of the second year of life. There is little strong evidence of typical age of first breeding but initial breeding by males could be inferred by the plumage of males in Utqiaġvik by plumage. At that stage, which the males were essentially all pure white, most were aged to about 3 to 4 years old. The snowy owl seems to markedly inconsistent in regard to breeding every year, often taking at least up to two years between attempts and sometimes as much as nearly a decade. 7 satellite-marking females in Canada proved that they did breed in consecutive years, with 1 breeding over 3 consecutive years. In 23 years at Utqiaġvik, snowys bred in 13 of them. Nesting success can reach 90–100% in even the largest clutches in high lemming years. While over the course of 21 years, 260 total nests were recorded in Utqiaġvik. There, from 4–54 nests were recorded annually. The Utqiaġvik nests bore 3 to 10 sized-clutches with a mean of 6 eggs per nest and an annual mean hatching success from 39 to 91%. 31–87% of chicks were able to depart on foot and 48–65% were annually estimated to survive to fledge; elsewhere, 40% survived to fledge. In another set, 97% of observed eggs both hatched and fledged. In Norway, the fledging success from 10 nests was much lower at about 46%. Norwegian data, which previously indicated it to be an almost accidental breeder in northern Norway, indicates that it is a more regular breeder than expected, though. 3 good years were found for snowy owls between 1968 and 2005: 1974 (when there were 12 pairs), 1978 (22 pairs) and 1985 (20 pairs), with 14 additional locations when potential (but not confirmed) breeding has occurred. The main determinable causes of nest failure were deemed to be starvation and exposure. A number of Norwegian and Finnish nests were known to fail due to severe black fly parasitism.
Longevity
The snowy owl can live a long life for a bird. Records show that the oldest snowy owls in captivity can live to 25 to even 30 years of age. Typical lifespans probably reach around 10 years in the wild. The longest known lifespan in the wild was of a snowy owl initially banded (possibly in its first winter) in Massachusetts and recovered dead in Montana 23 years and 10 months later. The annual survival rate for twelve females on Bylot Island was estimated at around 85–92.3%. It is often reputed that snowy owls frequently died from starvation, with historical accounts opining that they "had to" leave their breeding grounds due to lemming "crashes" but would starve to the south. However, it was proven fairly early on that snowy owls often do survive throughout the winter. This is reinforced somewhat by small radio-tracking and banding studies of snowy owls in the northern Great Plains and the intermountain valleys of the northwestern United States. More circumstantial evidence shows a lack of starvation in the eastern part of North America as well. There is evidence that some adults are known to return to the same wintering areas in ensuing years, areas which are far south of their breeding range. At Logan Airport, most snowy owls that are seen appear to be in good condition. Of 71 dead snowy owls found in winter in the northern Great Plains, 86% died from assorted traumas, including collisions with automobiles and other, usually manmade, objects as well as electrocutions and shootings. Only 14% of the 71 deaths were due to apparent starving. Data showed some owls appeared to incur injuries but healed and survived. More evidence was found in wintering snowy owls in New York of healed fractures, though some may require surgery to recover. 537 wintering birds in Saskatchewan were studied based on fat reserves, which were superior in females over males and adults over juveniles; while 31% of females lacked fat reserves, at least 45% of males found starving or in a state of infirmity were males and 63% turned into Wildlife rehabilitation centres were also males. In British Columbia, of 177 snowy owl deaths, of owls to die, only a small percentage were due to natural causes, such as assumed starvation at 13% and 12% were "found dead". One fledgling on Fetlar died due to pneumonia and Staphylococcus, while a second died from Aspergillosis. Evidence shows that in Utqiaġvik during exceptionally prolonged rains (i.e. 2 to 3 days), nest-departed young were vulnerable to starvation, leading to hypothermia and pneumonia. Due to their natural history, the snowy owl may be affected more severely by blood parasitism than other raptors, due to lowered immunity. Conversely, they appear to have lower levels of ectoparasites such as chewing lice than other large owls per large samples from Manitoba. The snowy owls averaged about 3.9 chewing lice per host against 7.5 for great grey owls and 10.5 for great horned owls.
Status
This species presence and numbers is dependent on amount of food available. In "lemming years", snowy owls can appear to be quite abundant in habitat. Numbers of snowy owls are difficult to estimate even within studies that take place over decades due to the nomadic nature of adults. The population of Scandinavia has long been perceived as very small and ephemeral with Finland holding 0–100 pairs; Norway holding 1–20 pairs and Sweden holding 1–50 pairs. A low breeding population within European Russia has been estimated to hold 1,300–4,500 pairs and Greenland to have 500–1,000 pairs. Other than northern part of the American continent, a majority of the snowy owl's breeding range is in northern Russia, but overall estimates are not known. An exact count of 4,871 individuals were seen on surveys between the Indigirka and Kolyma rivers. The numbers estimated by Partners in Flight and other authors by the 2000s was that North America held about 72,500 snowy owls, about 30% of which were juveniles. The Canadian population of snowy owls was estimated at 10,000–30,000 (in the 1990s) or even to 50,000–100,000 individuals, perhaps improbably. Within Canada, the population on Banks Island was once claimed at up to 15,000–25,000 in productive years and in Queen Elizabeth Islands at about 932 individuals. Alaska is the only state with breeding snowy owls but has probably quite a bit fewer breeding owls than does Canada. Furthermore, the Partners in Flight and the IUCN estimated that the world population was roughly 200,000–290,000 individuals as recently as the 2000s. However, in the 2010s, it has been discovered that all prior estimates were extremely excessive and that more precise numbers could be estimated with better surveying, phylogeographic data and more insights into the owl's free-wheeling wanderings. It is now believed that there are only 14,000–28,000 mature breeding pairs of snowy owls in the world. During lemming declines, the number of nesting females may drop down to as low as 1,700 worldwide, a dangerously low number, and the number of snowy owls worldwide is less than 10% of what it was once thought to be. Due to the small and rapidly declining population, the snowy was uplisted in 2017 to being a vulnerable species by the IUCN. A 52% decline has been inferred for the North American population since the 1960s with another even more drastic estimate placing the decline from 1970 to 2014 at 64%. Trends are harder to delineate in Scandinavia but a similar downward trend is thought to be occurring. Snowy owls are listed in Appendix II of the Convention on International Trade in Endangered Species (CITES) meaning international trade (including in parts and derivatives) is regulated.
Anthropogenic mortality and persecution
Of 438 band encounters in the USG banding laboratory, almost all causes of death that could be determined, whether intentional or not, were correlated with human interference. 34.2% or 150 were dead due to unknown causes, 11.9% were shot, 7.1% were hit by automobiles, 5.5% were found dead or injured on highways, 3.9% were collision from towers or wires, 2.7% were in animal traps, 2.1% in airplane birdstrikes, 0.6% were entangled while the remaining 33.3% recovered injured due to assorted or unknown causes. Snowy owls are endangered by heavy airport usage resulting in birdstrikes. Many such collisions are known in Canada and likely also in Siberia and Mongolia . Despite their danger to planes, no human fatalities have been recorded in collisions with this species. Snowy owls are always far outnumbered in Canadian airports in winter by short-eared owls. However, relative to its scarcity, the snowy accounts for a very large balance of the birdstrikes recorded at American airports due to the attractiveness of the habitat, accounting for 4.6% of 2456 recorded collisions (the barn owl is the most frequently involved in birdstrikes). The species is locally vulnerable to pesticides. The placement of buildings in the Utqiaġvik is now thought to have displaced some snowy owls. In Norway, potential sources of disturbance near the nests include tourism, recreation, reindeer husbandry, motorized traffic, dogs, photographers, ornithologists and scientists. Some biologist have expressed concern that radio-tagging of snowy owls may cause some unclear detrimental effect on snowy owls but little evidence is known if they actually make the owls more susceptible to death.
Snowy owls can be quite wary, as they are not infrequently hunted by Circumpolar peoples. Historically, the snowy owl was one of the most persecuted owl species. In the irruption of 1876–77, an estimated 500 snowy owls were shot, with similar numbers in 1889–90 and an estimated 500–1,000 killed in Ontario alone during 1901–02 invasion and about 800 killed in the 1905–06 invasion. Indigenous people of the Arctic historically killed snowy owls as food but now many communities in northern Alaska are fairly modernized, therefore biologists feel that the permitted killing of snowy owls by the indigenous is outdated. The consumption of snowy owls by humans has been proven as far back as ancient cave deposits in France and elsewhere, and they have even been considered as one of the most frequent food species for early humans. They do not shun developed areas especially with old field that hold rodents and, due to lack of human experience, can be extremely tame and unable to escape armed humans. In British Columbia, of 177 snowy owl deaths, the most often diagnosed cause of death was shootings at 25%, often well after legal protection of the species. The number poached snowy owls in Ontario is opined to be unusually high considering their scarcity. While the species was once otherwise killed as food and then later shot out of resentment for perceived threats against domestic and favored game stock, the reasoning behind ongoing shooting of snowy owls into the 21st century is not well-understood. Siberian snowy owls are frequently victim to baited fox traps, with possibly up to around 300 killed in a year based upon very rough estimates. Warfarin poisoning in use as rodenticides are known to kill some wintering snowy owls, including up to six at Logan Airport alone. Mercury concentrations, most likely through bioaccumulation, have been detected in snowy owls in the Aleutian Islands but it is not known whether fatal mercury poisoning has occurred. PCBs may have killed some snowy owls in concentration. Some airports have advocated and instituted the practice of shooting owls to avoid birdstrikes but successful translocation is possible and preferred given the species protected status.
Climate change is now widely perceived to perhaps the primary driver of the snowy owl's decline. As temperatures continue to rise, abiotic factors such as increased rain and reduced snow are likely to effect lemming populations and, in turn, snowy owls. These and potentially many other issues (possibly including modifying migrating behavior, vegetation composition, increased insect, disease and parasite activities, risk of hyperthermia) are a matter of concern. Additionally, reduction of sea ice, which snowy owls are now known to rely extensively on, as a result of warming climates, impacts could be significant. The effect of climate change was essentially confirmed in northern Greenland where a perhaps irrevocable collapse of the lemming population was observed. From 1998 to 2000, the lemming numbers appeared to have quickly declined. The number of lemmings per hectare (ha) is less than one-fifth of what it once was in Greenland (i.e. from 12 lemmings per ha to less than 2 per ha at peak). This is almost certainly correlated with a 98% decline in owl productivity as well as that of the local stoats (the long-tailed jaeger and Arctic foxes, though previously thought to be almost as reliant on lemmings, seem to be more loosely coupled and more generalized and did not decline as much). The amount of lemming mounds is much less than it once in northern Greenland and any variety of population cycle has been apparently abandoned by what remains of the lemmings.
In popular culture
The Harry Potter books by J. K. Rowling, and subsequent films of the same name, feature a female snowy owl named Hedwig. Concern was expressed by some in the media that the popularity of the Harry Potter films would cause an increase in the illicit owl trade of snowy owls. However, there was no strong evidence of an increase in snowy owls confiscated from the black market, despite a larger than typical number of snowy owls being reported at wildlife centres.
The snowy owl ( in French) is the avian symbol of Quebec.
| Biology and health sciences | Strigiformes | Animals |
252311 | https://en.wikipedia.org/wiki/Rule%20of%20inference | Rule of inference | In the philosophy of logic and logic, specifically in deductive reasoning, a rule of inference, inference rule or transformation rule is a logical form consisting of a function which takes premises, analyzes their syntax, and returns a conclusion (or conclusions).
For example, the rule of inference called modus ponens takes two premises, one in the form "If p then q" and another in the form "p", and returns the conclusion "q". The rule is valid with respect to the semantics of classical logic (as well as the semantics of many other non-classical logics), in the sense that if the premises are true (under an interpretation), then so is the conclusion.
Typically, a rule of inference preserves truth, a semantic property. In many-valued logic, it preserves a general designation. But a rule of inference's action is purely syntactic, and does not need to preserve any semantic property: any function from sets of formulae to formulae counts as a rule of inference. Usually only rules that are recursive are important; i.e. rules such that there is an effective procedure for determining whether any given formula is the conclusion of a given set of formulae according to the rule. An example of a rule that is not effective in this sense is the infinitary ω-rule.
Popular rules of inference in propositional logic include modus ponens, modus tollens, and contraposition. First-order predicate logic uses rules of inference to deal with logical quantifiers.
Standard form
In formal logic (and many related areas), rules of inference are usually given in the following standard form:
Premise#1
Premise#2
...
Premise#n
Conclusion
This expression states that whenever in the course of some logical derivation the given premises have been obtained, the specified conclusion can be taken for granted as well. The exact formal language that is used to describe both premises and conclusions depends on the actual context of the derivations. In a simple case, one may use logical formulae, such as in:
This is the modus ponens rule of propositional logic. Rules of inference are often formulated as schemata employing metavariables. In the rule (schema) above, the metavariables A and B can be instantiated to any element of the universe (or sometimes, by convention, a restricted subset such as propositions) to form an infinite set of inference rules.
A proof system is formed from a set of rules chained together to form proofs, also called derivations. Any derivation has only one final conclusion, which is the statement proved or derived. If premises are left unsatisfied in the derivation, then the derivation is a proof of a hypothetical statement: "if the premises hold, then the conclusion holds."
Example: Hilbert systems for two propositional logics
In a Hilbert system, the premises and conclusion of the inference rules are simply formulae of some language, usually employing metavariables. For graphical compactness of the presentation and to emphasize the distinction between axioms and rules of inference, this section uses the sequent notation () instead of a vertical presentation of rules.
In this notation,
is written as .
The formal language for classical propositional logic can be expressed using just negation (¬), implication (→) and propositional symbols. A well-known axiomatization, comprising three axiom schemata and one inference rule (modus ponens), is:
(CA1) ⊢ A → (B → A)
(CA2) ⊢ (A → (B → C)) → ((A → B) → (A → C))
(CA3) ⊢ (¬A → ¬B) → (B → A)
(MP) A, A → B ⊢ B
It may seem redundant to have two notions of inference in this case, ⊢ and →. In classical propositional logic, they indeed coincide; the deduction theorem states that A ⊢ B if and only if ⊢ A → B. There is however a distinction worth emphasizing even in this case: the first notation describes a deduction, that is an activity of passing from sentences to sentences, whereas A → B is simply a formula made with a logical connective, implication in this case. Without an inference rule (like modus ponens in this case), there is no deduction or inference. This point is illustrated in Lewis Carroll's dialogue called "What the Tortoise Said to Achilles", as well as later attempts by Bertrand Russell and Peter Winch to resolve the paradox introduced in the dialogue.
For some non-classical logics, the deduction theorem does not hold. For example, the three-valued logic of Łukasiewicz can be axiomatized as:
(CA1) ⊢ A → (B → A)
(LA2) ⊢ (A → B) → ((B → C) → (A → C))
(CA3) ⊢ (¬A → ¬B) → (B → A)
(LA4) ⊢ ((A → ¬A) → A) → A
(MP) A, A → B ⊢ B
This sequence differs from classical logic by the change in axiom 2 and the addition of axiom 4. The classical deduction theorem does not hold for this logic, however a modified form does hold, namely A ⊢ B if and only if ⊢ A → (A → B).
Admissibility and derivability
In a set of rules, an inference rule could be redundant in the sense that it is admissible or derivable. A derivable rule is one whose conclusion can be derived from its premises using the other rules. An admissible rule is one whose conclusion holds whenever the premises hold. All derivable rules are admissible. To appreciate the difference, consider the following set of rules for defining the natural numbers (the judgment asserts the fact that is a natural number):
The first rule states that 0 is a natural number, and the second states that s(n) is a natural number if n is. In this proof system, the following rule, demonstrating that the second successor of a natural number is also a natural number, is derivable:
Its derivation is the composition of two uses of the successor rule above. The following rule for asserting the existence of a predecessor for any nonzero number is merely admissible:
This is a true fact of natural numbers, as can be proven by induction. (To prove that this rule is admissible, assume a derivation of the premise and induct on it to produce a derivation of .) However, it is not derivable, because it depends on the structure of the derivation of the premise. Because of this, derivability is stable under additions to the proof system, whereas admissibility is not. To see the difference, suppose the following nonsense rule were added to the proof system:
In this new system, the double-successor rule is still derivable. However, the rule for finding the predecessor is no longer admissible, because there is no way to derive . The brittleness of admissibility comes from the way it is proved: since the proof can induct on the structure of the derivations of the premises, extensions to the system add new cases to this proof, which may no longer hold.
Admissible rules can be thought of as theorems of a proof system. For instance, in a sequent calculus where cut elimination holds, the cut rule is admissible.
| Mathematics | Mathematical logic | null |
252372 | https://en.wikipedia.org/wiki/Planetary%20system | Planetary system | A planetary system is a set of gravitationally bound non-stellar bodies in or out of orbit around a star or star system. Generally speaking, systems with one or more planets constitute a planetary system, although such systems may also consist of bodies such as dwarf planets, asteroids, natural satellites, meteoroids, comets, planetesimals and circumstellar disks. For example, the Sun together with the planetary system revolving around it, including Earth, form the Solar System. The term exoplanetary system is sometimes used in reference to other planetary systems.
Debris disks are known to be common while other objects are more difficult to observe.
Of particular interest to astrobiology is the habitable zone of planetary systems where planets could have surface liquid water, and thus, the capacity to support Earth-like life.
History
Heliocentrism
Heliocentrism is the doctrine that the Sun is at the centre of the universe, as opposed to geocentrism (placing Earth at the centre of the universe).
The notion of a heliocentric Solar System with the Sun at its centre is possibly first suggested in the Vedic literature of ancient India, which often refers to the Sun as the "centre of spheres". Some interpret Aryabhatta's writings in Āryabhaṭīya as implicitly heliocentric.
The idea was first proposed in Western philosophy and Greek astronomy as early as the 3rd century BC by Aristarchus of Samos, but received no support from most other ancient astronomers.
Discovery of the Solar System
De revolutionibus orbium coelestium by Nicolaus Copernicus, published in 1543, presented the first mathematically predictive heliocentric model of a planetary system. 17th-century successors Galileo Galilei, Johannes Kepler, and Sir Isaac Newton developed an understanding of physics which led to the gradual acceptance of the idea that the Earth moves around the Sun and that the planets are governed by the same physical laws that governed Earth.
Speculation on extrasolar planetary systems
In the 16th century the Italian philosopher Giordano Bruno, an early supporter of the Copernican theory that Earth and other planets orbit the Sun, put forward the view that the fixed stars are similar to the Sun and are likewise accompanied by planets. He was burned at the stake for his ideas by the Roman Inquisition.
In the 18th century, the same possibility was mentioned by Sir Isaac Newton in the "General Scholium" that concludes his Principia. Making a comparison to the Sun's planets, he wrote "And if the fixed stars are the centres of similar systems, they will all be constructed according to a similar design and subject to the dominion of One."
His theories gained popularity through the 19th and 20th centuries despite a lack of supporting evidence. Long before their confirmation by astronomers, conjecture on the nature of planetary systems had been a focus of the search for extraterrestrial intelligence and has been a prevalent theme in fiction, particularly science fiction.
Detection of exoplanets
The first confirmed detection of an exoplanet was in 1992, with the discovery of several terrestrial-mass planets orbiting the pulsar PSR B1257+12. The first confirmed detection of exoplanets of a main-sequence star was made in 1995, when a giant planet, 51 Pegasi b, was found in a four-day orbit around the nearby G-type star 51 Pegasi. The frequency of detections has increased since then, particularly through advancements in methods of detecting extrasolar planets and dedicated planet-finding programs such as the Kepler mission.
Origin and evolution
Planetary systems come from protoplanetary disks that form around stars as part of the process of star formation.
During formation of a system, much material is gravitationally-scattered into distant orbits, and some planets are ejected completely from the system, becoming rogue planets.
Evolved systems
High-mass stars
Planets orbiting pulsars have been discovered. Pulsars are the remnants of the supernova explosions of high-mass stars, but a planetary system that existed before the supernova would likely be mostly destroyed. Planets would either evaporate, be pushed off of their orbits by the masses of gas from the exploding star, or the sudden loss of most of the mass of the central star would see them escape the gravitational hold of the star, or in some cases the supernova would kick the pulsar itself out of the system at high velocity so any planets that had survived the explosion would be left behind as free-floating objects. Planets found around pulsars may have formed as a result of pre-existing stellar companions that were almost entirely evaporated by the supernova blast, leaving behind planet-sized bodies. Alternatively, planets may form in an accretion disk of fallback matter surrounding a pulsar. Fallback disks of matter that failed to escape orbit during a supernova may also form planets around black holes.
Lower-mass stars
As stars evolve and turn into red giants, asymptotic giant branch stars, and planetary nebulae they engulf the inner planets, evaporating or partially evaporating them depending on how massive they are. As the star loses mass, planets that are not engulfed move further out from the star.
If an evolved star is in a binary or multiple system, then the mass it loses can transfer to another star, forming new protoplanetary disks and second- and third-generation planets which may differ in composition from the original planets, which may also be affected by the mass transfer.
System architectures
The Solar System consists of an inner region of small rocky planets and outer region of large giant planets. However, other planetary systems can have quite different architectures. Studies suggest that architectures of planetary systems are dependent on the conditions of their initial formation. Many systems with a hot Jupiter gas giant very close to the star have been found. Theories, such as planetary migration or scattering, have been proposed for the formation of large planets close to their parent stars.
At present, few systems have been found to be analogous to the Solar System with terrestrial planets close to the parent star. More commonly, systems consisting of multiple Super-Earths have been detected.
Classification
Planetary system architectures may be partitioned into four classes based on how the mass of the planets is distributed around the host star:
Similar: The masses of all planets in a system are similar to each other. This architecture class is the most commonly-observed in our galaxy. Examples include Trappist-1. The planets in these systems are said to be like 'peas in a pod'.
Mixed: The masses of planets in a system show large increasing or decreasing variations. Examples of such systems are Gliese 876 and Kepler-89.
Anti-Ordered: The massive planets of a system are close to the star and smaller planets are further away from the star. There are currently no known examples of this architecture class.
Ordered: The mass of the planets in a system tends to increase with increasing distance from the host star. The Solar System, with small rocky planets in the inner part and giant planets in the outer part, is a type of Ordered system.
Peas in a pod
Multiplanetary systems tend to be in a "peas in a pod" configuration meaning they tend to have the following factors:
Size: planets within a system tend to be either similar or ordered in size.
Mass: planets within a system tend to be either similar or ordered in mass.
Spacing: planets within a system tend to be equally spaced apart.
Packing: small planets tend to be closely packed together, while large planets tend to have larger spacing.
Components
Planets and stars
Most known exoplanets orbit stars roughly similar to the Sun: that is, main-sequence stars of spectral categories F, G, or K. One reason is that planet-search programs have tended to concentrate on such stars. In addition, statistical analyses indicate that lower-mass stars (red dwarfs, of spectral category M) are less likely to have planets massive enough to be detected by the radial-velocity method. Nevertheless, several tens of planets around red dwarfs have been discovered by the Kepler space telescope by the transit method, which can detect smaller planets.
Circumstellar disks and dust structures
After planets, circumstellar disks are one of the most commonly-observed properties of planetary systems, particularly of young stars. The Solar System possesses at least four major circumstellar disks (the asteroid belt, Kuiper belt, scattered disc, and Oort cloud) and clearly-observable disks have been detected around nearby solar analogs including Epsilon Eridani and Tau Ceti. Based on observations of numerous similar disks, they are assumed to be quite common attributes of stars on the main sequence.
Interplanetary dust clouds have been studied in the Solar System and analogs are believed to be present in other planetary systems. Exozodiacal dust, an exoplanetary analog of zodiacal dust, the 1–100 micrometre-sized grains of amorphous carbon and silicate dust that fill the plane of the Solar System has been detected around the 51 Ophiuchi, Fomalhaut, Tau Ceti, and Vega systems.
Comets
there are 5,253 known Solar System comets and they are thought to be common components of planetary systems. The first exocomets were detected in 1987 around Beta Pictoris, a very young A-type main-sequence star. There are now a total of 11 stars around which the presence of exocomets have been observed or suspected. All discovered exocometary systems (Beta Pictoris, HR 10, 51 Ophiuchi, HR 2174, 49 Ceti, 5 Vulpeculae, 2 Andromedae, HD 21620, HD 42111, HD 110411, and more recently HD 172555) are around very young A-type stars.
Other components
Computer modelling of an impact in 2013 detected around the star NGC 2547-ID8 by the Spitzer Space Telescope, and confirmed by ground observations, suggests the involvement of large asteroids or protoplanets similar to the events believed to have led to the formation of terrestrial planets like the Earth.
Based on observations of the Solar System's large collection of natural satellites, they are believed common components of planetary systems; however, the existence of exomoons has not yet been confirmed. The star 1SWASP J140747.93-394542.6, in the constellation Centaurus, is a strong candidate for a natural satellite. Indications suggest that the confirmed extrasolar planet WASP-12b also has at least one satellite.
Orbital configurations
Unlike the Solar System, which has orbits that are nearly circular, many of the known planetary systems display much higher orbital eccentricity. An example of such a system is 16 Cygni.
Mutual inclination
The mutual inclination between two planets is the angle between their orbital planes. Many compact systems with multiple close-in planets interior to the equivalent orbit of Venus are expected to have very low mutual inclinations, so the system (at least the close-in part) would be even flatter than the Solar System. Captured planets could be captured into any arbitrary angle to the rest of the system. there are only a few systems where mutual inclinations have actually been measured One example is the Upsilon Andromedae system: the planets c and d have a mutual inclination of about 30 degrees.
Orbital dynamics
Planetary systems can be categorized according to their orbital dynamics as resonant, non-resonant-interacting, hierarchical, or some combination of these. In resonant systems the orbital periods of the planets are in integer ratios. The Kepler-223 system contains four planets in an 8:6:4:3 orbital resonance.
Giant planets are found in mean-motion resonances more often than smaller planets.
In interacting systems the planets' orbits are close enough together that they perturb the orbital parameters. The Solar System could be described as weakly interacting. In strongly interacting systems Kepler's laws do not hold.
In hierarchical systems the planets are arranged so that the system can be gravitationally considered as a nested system of two-bodies, e.g. in a star with a close-in hot Jupiter with another gas giant much further out, the star and hot Jupiter form a pair that appears as a single object to another planet that is far enough out.
Other, as yet unobserved, orbital possibilities include: double planets; various co-orbital planets such as quasi-satellites, trojans and exchange orbits; and interlocking orbits maintained by precessing orbital planes.
Number of planets, relative parameters and spacings
On The Relative Sizes of Planets Within Kepler Multiple Candidate Systems, David R. Ciardi et al. December 9, 2012
The Kepler Dichotomy among the M Dwarfs: Half of Systems Contain Five or More Coplanar Planets, Sarah Ballard, John Asher Johnson, October 15, 2014
Exoplanet Predictions Based on the Generalised Titius-Bode Relation, Timothy Bovaird, Charles H. Lineweaver, August 1, 2013
The Solar System and the Exoplanet Orbital Eccentricity - Multiplicity Relation, Mary Anne Limbach, Edwin L. Turner, April 9, 2014
The period ratio distribution of Kepler's candidate multiplanet systems, Jason H. Steffen, Jason A. Hwang, September 11, 2014
Are Planetary Systems Filled to Capacity? A Study Based on Kepler Results, Julia Fang, Jean-Luc Margot, February 28, 2013
Planet capture
Free-floating planets in open clusters have similar velocities to the stars and so can be recaptured. They are typically captured into wide orbits between 100 and 105 AU. The capture efficiency decreases with increasing cluster size, and for a given cluster size it increases with the host/primary mass. It is almost independent of the planetary mass. Single and multiple planets could be captured into arbitrary unaligned orbits, non-coplanar with each other or with the stellar host spin, or pre-existing planetary system. Some planet–host metallicity correlation may still exist due to the common origin of the stars from the same cluster. Planets would be unlikely to be captured around neutron stars because these are likely to be ejected from the cluster by a pulsar kick when they form. Planets could even be captured around other planets to form free-floating planet binaries. After the cluster has dispersed some of the captured planets with orbits larger than 106 AU would be slowly disrupted by the galactic tide and likely become free-floating again through encounters with other field stars or giant molecular clouds.
Zones
Habitable zone
The habitable zone around a star is the region where the temperature range allows for liquid water to exist on a planet; that is, not too close to the star for the water to evaporate and not too far away from the star for the water to freeze. The heat produced by stars varies depending on the size and age of the star; this means the habitable zone will also vary accordingly. Also, the atmospheric conditions on the planet influence the planet's ability to retain heat so that the location of the habitable zone is also specific to each type of planet.
Habitable zones have usually been defined in terms of surface temperature; however, over half of Earth's biomass is from subsurface microbes, and temperature increases as depth underground increases, so the subsurface can be conducive for life when the surface is frozen; if this is considered, the habitable zone extends much further from the star.
Studies in 2013 indicate that an estimated 22±8% of Sun-like stars have an Earth-sized planet in the habitable zone.
Venus zone
The Venus zone is the region around a star where a terrestrial planet would have runaway greenhouse conditions like Venus, but not so near the star that the atmosphere completely escapes. As with the habitable zone, the location of the Venus zone depends on several factors, including the type of star and properties of the planets such as mass, rotation rate, and atmospheric clouds. Studies of the Kepler spacecraft data indicate that 32% of red dwarfs have potentially Venus-like planets based on planet size and distance from star, increasing to 45% for K-type and G-type stars. Several candidates have been identified, but spectroscopic follow-up studies of their atmospheres are required to determine whether they are like Venus.
Galactic distribution of planets
The Milky Way is 100,000 light-years across, but 90% of planets with known distances are within about 2000 light years of Earth, as of July 2014. One method that can detect planets much further away is microlensing. The upcoming Nancy Grace Roman Space Telescope could use microlensing to measure the relative frequency of planets in the galactic bulge versus the galactic disk. So far, the indications are that planets are more common in the disk than the bulge. Estimates of the distance of microlensing events is difficult: the first planet considered with high probability of being in the bulge is MOA-2011-BLG-293Lb at a distance of 7.7 kiloparsecs (about 25,000 light years).
Population I, or metal-rich stars, are those young stars whose metallicity is highest. The high metallicity of population I stars makes them more likely to possess planetary systems than older populations, because planets form by the accretion of metals. The Sun is an example of a metal-rich star. These are common in the spiral arms of the Milky Way. Generally, the youngest stars, the extreme population I, are found farther in and intermediate population I stars are farther out, etc. The Sun is considered an intermediate population I star. Population I stars have regular elliptical orbits around the Galactic Center, with a low relative velocity.
Population II, or metal-poor stars, are those with relatively low metallicity which can have hundreds (e.g. BD +17° 3248) or thousands (e.g. Sneden's Star) times less metallicity than the Sun. These objects formed during an earlier time of the universe. Intermediate population II stars are common in the bulge near the center of the Milky Way, whereas Population II stars found in the galactic halo are older and thus more metal-poor. Globular clusters also contain high numbers of population II stars.
In 2014, the first planets around a halo star were announced around Kapteyn's star, the nearest halo star to Earth, around 13 light years away. However, later research suggests that Kapteyn b is just an artefact of stellar activity and that Kapteyn c needs more study to be confirmed. The metallicity of Kapteyn's star is estimated to be about 8 times less than the Sun.
Different types of galaxies have different histories of star formation and hence planet formation. Planet formation is affected by the ages, metallicities, and orbits of stellar populations within a galaxy. Distribution of stellar populations within a galaxy varies between the different types of galaxies.
Stars in elliptical galaxies are much older than stars in spiral galaxies. Most elliptical galaxies contain mainly low-mass stars, with minimal star-formation activity. The distribution of the different types of galaxies in the universe depends on their location within galaxy clusters, with elliptical galaxies found mostly close to their centers.
| Physical sciences | Stellar astronomy | null |
252563 | https://en.wikipedia.org/wiki/Recreational%20vehicle | Recreational vehicle | A recreational vehicle, often abbreviated as RV, is a motor vehicle or trailer that includes living quarters designed for accommodation. Types of RVs include motorhomes, campervans, coaches, caravans (also known as travel trailers and campers), fifth-wheel trailers, popup campers, and truck campers.
Typical amenities of an RV include a kitchen, a bathroom, and one or more beds. RVs can range from utilitarian – containing only sleeping quarters and basic cooking facilities – to luxurious, with features like air conditioning (AC), water heaters, televisions and satellite receivers, and quartz countertops.
Types
RVs can be either trailers that are towed by vehicles or vehicles that drive themselves. Most RVs have one level, but there are also some with two levels. To save space while traveling, larger RVs often have slide-outs or canopies that open up when parked. Self-driving RVs can be categorized into pushers, which have the engine at the back, and pullers, which have the engine at the front. Generally, pushers use diesel fuel, while pullers use gasoline.
History
The first recreational vehicles were horse-drawn. They evolved during the second half of the nineteenth century as adaptations of vehicles used for other purposes, including public transport caravans (UK, also known as stage wagons), gypsy vardos (Europe), living vans (UK), ambulance wagons (US) and sheep herders wagons (US).
The first, currently-known, purpose-built RV was the horse-drawn Wanderer (UK), commissioned from the Bristol Wagon Works Company by Dr. Gordon Stables in 1884. Stables was a pioneer of the UK's Gentlemen Gypsy movement (1885–1914) which promoted the restorative benefits of horse-drawn leisure caravanning and inspired the formation of the world's first RV club, The Caravan Club (UK), in 1907. The Wanderer was closely followed by the McMaster Camping Car (US,1889).Camping-vehicle In the 1890s, US RV pioneers self-built timber 'houses on wheels' for health, leisure and hunting purposes. The most widely reported of these were those of Morgan Lasley and his family. Horse-drawn RV use declined after the First World War as many horses were killed during the war and automobiles became cheaper, more powerful and more widely available.
The first powered RVs were steam-drawn trailers from France including the Grande Diligence of Prince Oldenburg (1896) and the De Dion Bouton trailer of Monsieur Rénodier (1898). The first steam-driven motorhome was the Quo Vadis (France,1900) and the first gasoline-driven motorhome was the Passe Partout (France, 1902). The first recorded powered motorhomes in America were the 'camp cars' of Roy Faye and Freeman Young of 1904–06 (a 1904 Rambler, 1905 Thomas Flyer and 1906 Matheson). Lightweight tent trailers were especially popular in the US from 1911, thanks to improved roads, new national parks and the affordability of tow vehicles such as the Ford Model T. At the other end of the price scale, luxury touring limousines, developed in France by De Dietrich in 1904, were built in small numbers in the US by Welch (1909) and Pierce Arrow (1910). The first US RV club, the Tin Can Tourists, was formed in 1919. The first known, recreational fifth wheeler was the Auto Salon Deluxe built in Belgium for Baron Crawhez by Auto-Mixte Pescatore in 1913.
In the early twentieth century RV builders in the UK (Navarac, Piggott Bros, Eccles, Bertram Hutchings), the US (Detroit Trailer Company, Welch, Graham Brothers, Pierce-Arrow) and France (De Dietrich, Cadel) experimented with a wide range of RV types including caravans and trailers, motorhomes, touring limousines, tent trailers and fifth wheelers. Early motorhomes ('house cars' in the US) were usually converted goods trucks and were heavy, noisy, inflexible and expensive, restricting their use to the wealthy or self-builders. The 'one box' RV was not seen in large numbers until the small, lightweight Volkswagen Kombi of 1950. During the 1920s and 1930s, caravans (travel trailers) became the dominant form of RV in the UK due to their low cost, weather-resistance and flexibility. There was likewise a travel trailer boom in the US in the 1930s as automobile production-line manufacturing techniques were used in travel trailer manufacturing to meet growing demand from recreational users and those seeking low-cost housing during The Great Depression. Self-built trailers were highly popular in the US during the 1930s and travel trailers featured in a number of Hollywood movies including Mickey's Trailer (1938).
The 1920s and 1930s, saw some influential, maverick builders construct innovative RVs in small numbers. These included Bertram Hutchings (UK, 1930–39, streamlined caravans), Charles Louvet (France, 1924–34, aircraft-inspired, coach-built motorhomes and trailers), Noel Pemberton Billing (UK, 1927, Road Yacht motorhome), Glenn Curtiss (US, 1918–30, Adams Motorbungalo, Curtiss Aerocar, Aero Coupler hitch) and William Hawley Bowlus (US, 1934, aluminum monocoque trailers). Wally Byam's Airstream (US, 1931 onwards) was a successful trailer builder of the period and is the only survivor of over 400 pre-WW2 US RV manufacturers. These early advancements in RV and trailer design established the foundation for a burgeoning industry that would continue to develop over the subsequent decades.
Prior to WW2 a number of other countries developed their own small-scale RV manufacturing industries including Germany, Australia, New Zealand and the Netherlands. Germany had a particular focus on small, lightweight caravans towable by low-cost automobiles whilst Australia developed 'pop-top' caravans with good ventilation, ample water storage and high ground clearance. During WW2 RV production was halted in most countries except when required as accommodation for military personnel or essential workers.
Between 1945 and 1960, RVs flourished in many western countries as disposable income and leisure time grew. Dedicated RV parks were established to cater to the needs of both short and long-term RV users. Improvements in RV technology including batteries, fridges, gas cookers, toilets and lightweight construction techniques dramatically improved RV comfort levels. More powerful gasoline and diesel engines allowed RVs to increase in size, weight and speed. Regulations were introduced in many countries to control how RVs were manufactured and used.
Industry
In the United States, about 85 percent of recreational vehicles sold are manufactured in Indiana, and roughly two-thirds of that production in Elkhart County, which calls itself "the RV Capital of the World", population 206,000. The industry has US$32.4 billion annual economic impact in Indiana, pays US$3.1 billion in taxes to the state and supports 126,140 jobs and US$7.8 billion in wages, according to the RV Industry Association.
The recreational vehicle industry around Elkhart is part of a large network of related transport equipment companies, including utility trailer makers and specialty bus manufacturers, who source from the same supply chains. The industry has taken hits from US tariffs on steel and aluminum and other duties on RV parts made in China, from plumbing fixtures to electronic components to vinyl seat covers. Tariff-related price hikes forced manufacturers to pass on some of the increased costs through higher RV prices, which in turn has contributed to slower sales. Shipments of RVs to dealers fell 22% percent in the first five months of 2019, compared to the same period a year earlier, after dropping 4% in 2018.
Usage
RVs are most commonly used for living quarters while traveling. People may choose to take a road trip in their RV and use the RV to sleep in, rather than a hotel room. They may even decide to tow their car from the back of the RV so they can use that to travel around more easily when they reach their destination.
Although the most common usage of an RV is as temporary accommodation when traveling, some people use an RV as their main residence. In fact, one million Americans live in RVs. In the United States and Canada, traveling south each winter to a warmer climate is referred to as snowbirding. In Australia, the slang term for a retired person who travels in a recreational vehicle is a "grey nomad".
There are local and national RV rental companies, such as Adventure KT and Outdoorsy that specialize in renting RVs to families for vacationing purposes. People enjoy the road trip and luxuries an RV provides while traveling without having a long-term expense. This is similar to home vacation rentals but is cheaper and also offers the flexibility of itinerary planning.
While it is legal in all of the United States to live in an RV, there are laws regarding where and for how long RVs can be parked.
Some owners fit solar panels to the roof of their RV. It is possible for RV users to live off the electrical grid while still having access to internet, making remote working feasible.
Usage of RVs is common at rural festivals such as Burning Man, but most festivals have strict rules about operating an RV during the event. Burning Man is strict about RV water leaks, and generator usage is another of the restrictions that festivals put on the use of RVs. Bluegrass Festivals regularly host RVs; they become the locations for afterhours jamming by participants.
Recently, RVs have been proposed as a partial solution to the homelessness problems found in cities across the US. RVs for Homeless (https://rvforhomeless.com/) has been formed as a coalition to encourage the donation of older vehicles to alleviate homelessness. This is distinct from the simple use of older vehicles as temporary shelters in homeless encampments.
Demographics
United States
As of 2016, the average age of RV owners in the United States was 45, a three-year decrease since 2015. Per 2020 research reports, more millennials are interested in buying RVs due to their increased demand for camping and outdoor recreational activities, especially in the US.
Due to the COVID-19 pandemic, sales of RVs in the United States have increased, and as of March 2021, 11.3 million households own an RV, which is a 26 percent increase over the past ten years. In the month of October 2021, 58,000 RVs were manufactured in North America, the most ever in a single month.
Current trends
Wholesale recreational vehicle (RV) shipments during the first four months of 2021 rose more than 86% over 2020, because consumer demand for RVs has soared since the pandemic's onset.
Analysts expect industry revenue to increase 1.1% through 2026 as U.S. residents continue to plan domestic trips because of continued travel restrictions around the globe.
| Technology | Motorized road transport | null |
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