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18955875 | https://en.wikipedia.org/wiki/Tree | Tree | In botany, a tree is a perennial plant with an elongated stem, or trunk, usually supporting branches and leaves. In some usages, the definition of a tree may be narrower, including only woody plants with secondary growth, plants that are usable as lumber or plants above a specified height. In wider definitions, the taller palms, tree ferns, bananas, and bamboos are also trees.
Trees are not a monophyletic taxonomic group but consist of a wide variety of plant species that have independently evolved a trunk and branches as a way to tower above other plants to compete for sunlight. The majority of tree species are angiosperms or hardwoods; of the rest, many are gymnosperms or softwoods. Trees tend to be long-lived, some trees reaching several thousand years old. Trees evolved around 370 million years ago, and it is estimated that there are around three trillion mature trees in the world currently.
A tree typically has many secondary branches supported clear of the ground by the trunk, which typically contains woody tissue for strength, and vascular tissue to carry materials from one part of the tree to another. For most trees the trunk is surrounded by a layer of bark which serves as a protective barrier. Below the ground, the roots branch and spread out widely; they serve to anchor the tree and extract moisture and nutrients from the soil. Above ground, the branches divide into smaller branches and shoots. The shoots typically bear leaves, which capture light energy and convert it into sugars by photosynthesis, providing the food for the tree's growth and development.
Trees usually reproduce using seeds. Flowering plants have their seeds inside fruits, while conifers carry their seeds in cones, and tree ferns produce spores instead.
Trees play a significant role in reducing erosion and moderating the climate. They remove carbon dioxide from the atmosphere and store large quantities of carbon in their tissues. Trees and forests provide a habitat for many species of animals and plants. Tropical rainforests are among the most biodiverse habitats in the world. Trees provide shade and shelter, timber for construction, fuel for cooking and heating, and fruit for food as well as having many other uses. In much of the world, forests are shrinking as trees are cleared to increase the amount of land available for agriculture. Because of their longevity and usefulness, trees have always been revered, with sacred groves in various cultures, and they play a role in many of the world's mythologies.
Definition
Although "tree" is a common word, there is no universally recognised precise definition of what a tree is, either botanically or in common language. In its broadest sense, a tree is any plant with the general form of an elongated stem, or trunk, which supports the photosynthetic leaves or branches at some distance above the ground. Trees are also typically defined by height, with smaller plants from being called shrubs, so the minimum height of a tree is only loosely defined. Large herbaceous plants such as papaya and bananas are trees in this broad sense.
A commonly applied narrower definition is that a tree has a woody trunk formed by secondary growth, meaning that the trunk thickens each year by growing outwards, in addition to the primary upwards growth from the growing tip. Under such a definition, herbaceous plants such as palms, bananas and papayas are not considered trees regardless of their height, growth form or stem girth. Certain monocots may be considered trees under a slightly looser definition; while the Joshua tree, bamboos and palms do not have secondary growth and never produce true wood with growth rings, they may produce "pseudo-wood" by lignifying cells formed by primary growth. Tree species in the genus Dracaena, despite also being monocots, do have secondary growth caused by meristem in their trunk, but it is different from the thickening meristem found in dicotyledonous trees.
Aside from structural definitions, trees are commonly defined by use; for instance, as those plants which yield lumber.
Overview
The tree growth habit is an evolutionary adaptation found in different groups of plants: by growing taller, trees are able to compete better for sunlight. Trees tend to be tall and long-lived, some reaching several thousand years old. Several trees are among the oldest organisms now living. Trees have modified structures such as thicker stems composed of specialised cells that add structural strength and durability, allowing them to grow taller than many other plants and to spread out their foliage. They differ from shrubs, which have a similar growth form, by usually growing larger and having a single main stem; but there is no consistent distinction between a tree and a shrub, made more confusing by the fact that trees may be reduced in size under harsher environmental conditions such as on mountains and subarctic areas. The tree form has evolved separately in unrelated classes of plants in response to similar environmental challenges, making it a classic example of parallel evolution. With an estimated 60,000-100,000 species, the number of trees worldwide might total twenty-five per cent of all living plant species. The greatest number of these grow in tropical regions; many of these areas have not yet been fully surveyed by botanists, making tree diversity and ranges poorly known.
The majority of tree species are angiosperms or hardwoods. Of the rest, many are gymnosperms or softwood trees; these include conifers, cycads, ginkgophytes and gnetales, which produce seeds which are not enclosed in fruits, but in open structures such as pine cones, and many have tough waxy leaves, such as pine needles. Most angiosperm trees are eudicots, the "true dicotyledons", so named because the seeds contain two cotyledons or seed leaves. There are also some trees among the old lineages of flowering plants called basal angiosperms or paleodicots; these include Amborella, Magnolia, nutmeg and avocado, while trees such as bamboo, palms and bananas are monocots.
Wood gives structural strength to the trunk of most types of tree; this supports the plant as it grows larger. The vascular system of trees allows water, nutrients and other chemicals to be distributed around the plant, and without it trees would not be able to grow as large as they do. Trees need to draw water high up the stem through the xylem from the roots by capillary action, as water continually evaporates from the leaves in the process of transpiration. If insufficient water is available the leaves will die. The three main parts of trees include the root, stem, and leaves; they are integral parts of the vascular system which interconnects all the living cells. In trees and other plants that develop wood, the vascular cambium allows the expansion of vascular tissue that produces woody growth. Because this growth ruptures the epidermis of the stem, woody plants also have a cork cambium that develops among the phloem. The cork cambium gives rise to thickened cork cells to protect the surface of the plant and reduce water loss. Both the production of wood and the production of cork are forms of secondary growth.
Trees are either evergreen, having foliage that persists and remains green throughout the year, or deciduous, shedding their leaves at the end of the growing season and then having a dormant period without foliage. Most conifers are evergreens, but larches (Larix and Pseudolarix) are deciduous, dropping their needles each autumn, and some species of cypress (Glyptostrobus, Metasequoia and Taxodium) shed small leafy shoots annually in a process known as cladoptosis. The crown is the spreading top of a tree including the branches and leaves, while the uppermost layer in a forest, formed by the crowns of the trees, is known as the canopy. A sapling is a young tree.
Many tall palms are herbaceous monocots, which do not undergo secondary growth and never produce wood. In many tall palms, the terminal bud on the main stem is the only one to develop, so they have unbranched trunks with large spirally arranged leaves. Some of the tree ferns, order Cyatheales, have tall straight trunks, growing up to , but these are composed not of wood but of rhizomes which grow vertically and are covered by numerous adventitious roots.
Distribution
The number of trees in the world, according to a 2015 estimate, is 3.04 trillion, of which 1.39 trillion (46%) are in the tropics or sub-tropics, 0.61 trillion (20%) in the temperate zones, and 0.74 trillion (24%) in the coniferous boreal forests. The estimate is about eight times higher than previous estimates, and is based on tree densities measured on over 400,000 plots. It remains subject to a wide margin of error, not least because the samples are mainly from Europe and North America. The estimate suggests that about 15 billion trees are cut down annually and about 5 billion are planted. In the 12,000 years since the start of human agriculture, the number of trees worldwide has decreased by 46%. There are approximately 64,100 known tree species in the world. With 43% of all tree species, South America has the highest biodiversity, followed by Eurasia (22%), Africa (16%), North America (15%), and Oceania (11%).
In suitable environments, such as the Daintree Rainforest in Queensland, or the mixed podocarp and broadleaf forest of Ulva Island, New Zealand, forest is the more-or-less stable climatic climax community at the end of a plant succession, where open areas such as grassland are colonised by taller plants, which in turn give way to trees that eventually form a forest canopy.
In cool temperate regions, conifers often predominate; a widely distributed climax community in the far north of the northern hemisphere is moist taiga or northern coniferous forest (also called boreal forest). Taiga is the world's largest land biome, forming 29% of the world's forest cover. The long cold winter of the far north is unsuitable for plant growth and trees must grow rapidly in the short summer season when the temperature rises and the days are long. Light is very limited under their dense cover and there may be little plant life on the forest floor, although fungi may abound. Similar woodland is found on mountains where the altitude causes the average temperature to be lower thus reducing the length of the growing season.
Where rainfall is relatively evenly spread across the seasons in temperate regions, temperate broadleaf and mixed forest typified by species like oak, beech, birch and maple is found. Temperate forest is also found in the southern hemisphere, as for example in the Eastern Australia temperate forest, characterised by Eucalyptus forest and open acacia woodland.
In tropical regions with a monsoon or monsoon-like climate, where a drier part of the year alternates with a wet period as in the Amazon rainforest, different species of broad-leaved trees dominate the forest, some of them being deciduous. In tropical regions with a drier savanna climate and insufficient rainfall to support dense forests, the canopy is not closed, and plenty of sunshine reaches the ground which is covered with grass and scrub. Acacia and baobab are well adapted to living in such areas.
Parts
Roots
The roots of a tree serve to anchor it to the ground and gather water and nutrients to transfer to all parts of the tree. They are also used for reproduction, defence, survival, energy storage and many other purposes. The radicle or embryonic root is the first part of a seedling to emerge from the seed during the process of germination. This develops into a taproot which goes straight downwards. Within a few weeks lateral roots branch out of the side of this and grow horizontally through the upper layers of the soil. In most trees, the taproot eventually withers away and the wide-spreading laterals remain. Near the tip of the finer roots are single cell root hairs. These are in immediate contact with the soil particles and can absorb water and nutrients such as potassium in solution. The roots require oxygen to respire and only a few species such as mangroves and the pond cypress (Taxodium ascendens) can live in permanently waterlogged soil.
In the soil, the roots encounter the hyphae of fungi. Many of these are known as mycorrhiza and form a mutualistic relationship with the tree roots. Some are specific to a single tree species, which will not flourish in the absence of its mycorrhizal associate. Others are generalists and associate with many species. The tree acquires minerals such as phosphorus from the fungus, while the fungus obtains the carbohydrate products of photosynthesis from the tree. The hyphae of the fungus can link different trees and a network is formed, transferring nutrients and signals from one place to another. The fungus promotes growth of the roots and helps protect the trees against predators and pathogens. It can also limit damage done to a tree by pollution as the fungus accumulate heavy metals within its tissues. Fossil evidence shows that roots have been associated with mycorrhizal fungi since the early Paleozoic, four hundred million years ago, when the first vascular plants colonised dry land.
Some trees such as Alder (Alnus species) have a symbiotic relationship with Frankia species, a filamentous bacterium that can fix nitrogen from the air, converting it into ammonia. They have actinorhizal root nodules on their roots in which the bacteria live. This process enables the tree to live in low nitrogen habitats where they would otherwise be unable to thrive. The plant hormones called cytokinins initiate root nodule formation, in a process closely related to mycorrhizal association.
It has been demonstrated that some trees are interconnected through their root system, forming a colony. The interconnections are made by the inosculation process, a kind of natural grafting or welding of vegetal tissues. The tests to demonstrate this networking are performed by injecting chemicals, sometimes radioactive, into a tree, and then checking for its presence in neighbouring trees.
The roots are, generally, an underground part of the tree, but some tree species have evolved roots that are aerial. The common purposes for aerial roots may be of two kinds, to contribute to the mechanical stability of the tree, and to obtain oxygen from air. An instance of mechanical stability enhancement is the red mangrove that develops prop roots that loop out of the trunk and branches and descend vertically into the mud. A similar structure is developed by the Indian banyan. Many large trees have buttress roots which flare out from the lower part of the trunk. These brace the tree rather like angle brackets and provide stability, reducing sway in high winds. They are particularly prevalent in tropical rainforests where the soil is poor and the roots are close to the surface.
Some tree species have developed root extensions that pop out of soil, in order to get oxygen, when it is not available in the soil because of excess water. These root extensions are called pneumatophores, and are present, among others, in black mangrove and pond cypress.
Trunk
The main purpose of the trunk is to raise the leaves above the ground, enabling the tree to overtop other plants and outcompete them for light. It also transports water and nutrients from the roots to the aerial parts of the tree, and distributes the food produced by the leaves to all other parts, including the roots.
In the case of angiosperms and gymnosperms, the outermost layer of the trunk is the bark, mostly composed of dead cells of phellem (cork). It provides a thick, waterproof covering to the living inner tissue. It protects the trunk against the elements, disease, animal attack and fire. It is perforated by a large number of fine breathing pores called lenticels, through which oxygen diffuses. Bark is continually replaced by a living layer of cells called the cork cambium or phellogen. The London plane (Platanus × hispanica) periodically sheds its bark in large flakes. Similarly, the bark of the silver birch (Betula pendula) peels off in strips. As the tree's girth expands, newer layers of bark are larger in circumference, and the older layers develop fissures in many species. In some trees such as the pine (Pinus species) the bark exudes sticky resin which deters attackers whereas in rubber trees (Hevea brasiliensis) it is a milky latex that oozes out. The quinine bark tree (Cinchona officinalis) contains bitter substances to make the bark unpalatable. Large tree-like plants with lignified trunks in the Pteridophyta, Arecales, Cycadophyta and Poales such as the tree ferns, palms, cycads and bamboos have different structures and outer coverings.
Although the bark functions as a protective barrier, it is itself attacked by boring insects such as beetles. These lay their eggs in crevices and the larvae chew their way through the cellulose tissues leaving a gallery of tunnels. This may allow fungal spores to gain admittance and attack the tree. Dutch elm disease is caused by a fungus (Ophiostoma species) carried from one elm tree to another by various beetles. The tree reacts to the growth of the fungus by blocking off the xylem tissue carrying sap upwards and the branch above, and eventually the whole tree, is deprived of nourishment and dies. In Britain in the 1990s, 25 million elm trees were killed by this disease.
The innermost layer of bark is known as the phloem and this is involved in the transport of the sap containing the sugars made by photosynthesis to other parts of the tree. It is a soft spongy layer of living cells, some of which are arranged end to end to form tubes. These are supported by parenchyma cells which provide padding and include fibres for strengthening the tissue. Inside the phloem is a layer of undifferentiated cells one cell thick called the vascular cambium layer. The cells are continually dividing, creating phloem cells on the outside and wood cells known as xylem on the inside.
The newly created xylem is the sapwood. It is composed of water-conducting cells and associated cells which are often living, and is usually pale in colour. It transports water and minerals from the roots to the upper parts of the tree. The oldest, inner part of the sapwood is progressively converted into heartwood as new sapwood is formed at the cambium. The conductive cells of the heartwood are blocked in some species. Heartwood is usually darker in colour than the sapwood. It is the dense central core of the trunk giving it rigidity. Three quarters of the dry mass of the xylem is cellulose, a polysaccharide, and most of the remainder is lignin, a complex polymer. A transverse section through a tree trunk or a horizontal core will show concentric circles of lighter or darker wood – tree rings. These rings are the annual growth rings There may also be rays running at right angles to growth rings. These are vascular rays which are thin sheets of living tissue permeating the wood. Many older trees may become hollow but may still stand upright for many years.
Buds and growth
Trees do not usually grow continuously throughout the year but mostly have spurts of active expansion followed by periods of rest. This pattern of growth is related to climatic conditions; growth normally ceases when conditions are either too cold or too dry. In readiness for the inactive period, trees form buds to protect the meristem, the zone of active growth. Before the period of dormancy, the last few leaves produced at the tip of a twig form scales. These are thick, small and closely wrapped and enclose the growing point in a waterproof sheath. Inside this bud there is a rudimentary stalk and neatly folded miniature leaves, ready to expand when the next growing season arrives. Buds also form in the axils of the leaves ready to produce new side shoots. A few trees, such as the eucalyptus, have "naked buds" with no protective scales and some conifers, such as the Lawson's cypress, have no buds but instead have little pockets of meristem concealed among the scale-like leaves.
When growing conditions improve, such as the arrival of warmer weather and the longer days associated with spring in temperate regions, growth starts again. The expanding shoot pushes its way out, shedding the scales in the process. These leave behind scars on the surface of the twig. The whole year's growth may take place in just a few weeks. The new stem is unlignified at first and may be green and downy. The Arecaceae (palms) have their leaves spirally arranged on an unbranched trunk. In some tree species in temperate climates, a second spurt of growth, a Lammas growth may occur which is believed to be a strategy to compensate for loss of early foliage to insect predators.
Primary growth is the elongation of the stems and roots. Secondary growth consists of a progressive thickening and strengthening of the tissues as the outer layer of the epidermis is converted into bark and the cambium layer creates new phloem and xylem cells. The bark is inelastic. Eventually the growth of a tree slows down and stops and it gets no taller. If damage occurs the tree may in time become hollow.
Leaves
Leaves are structures specialised for photosynthesis and are arranged on the tree in such a way as to maximise their exposure to light without shading each other. They are an important investment by the tree and may be thorny or contain phytoliths, lignins, tannins or poisons to discourage herbivory. Trees have evolved leaves in a wide range of shapes and sizes, in response to environmental pressures including climate and predation. They can be broad or needle-like, simple or compound, lobed or entire, smooth or hairy, delicate or tough, deciduous or evergreen. The needles of coniferous trees are compact but are structurally similar to those of broad-leaved trees. They are adapted for life in environments where resources are low or water is scarce. Frozen ground may limit water availability and conifers are often found in colder places at higher altitudes and higher latitudes than broad leaved trees. In conifers such as fir trees, the branches hang down at an angle to the trunk, enabling them to shed snow. In contrast, broad leaved trees in temperate regions deal with winter weather by shedding their leaves. When the days get shorter and the temperature begins to decrease, the leaves no longer make new chlorophyll and the red and yellow pigments already present in the blades become apparent. Synthesis in the leaf of a plant hormone called auxin also ceases. This causes the cells at the junction of the petiole and the twig to weaken until the joint breaks and the leaf floats to the ground. In tropical and subtropical regions, many trees keep their leaves all year round. Individual leaves may fall intermittently and be replaced by new growth but most leaves remain intact for some time. Other tropical species and those in arid regions may shed all their leaves annually, such as at the start of the dry season. Many deciduous trees flower before the new leaves emerge. A few trees do not have true leaves but instead have structures with similar external appearance such as Phylloclades – modified stem structures – as seen in the genus Phyllocladus.
Reproduction
Trees can be pollinated either by wind or by animals, mostly insects. Many angiosperm trees are insect pollinated. Wind pollination may take advantage of increased wind speeds high above the ground. Trees use a variety of methods of seed dispersal. Some rely on wind, with winged or plumed seeds. Others rely on animals, for example with edible fruits. Others again eject their seeds (ballistic dispersal), or use gravity so that seeds fall and sometimes roll.
Seeds
Seeds are the primary way that trees reproduce and their seeds vary greatly in size and shape. Some of the largest seeds come from trees, but the largest tree, Sequoiadendron giganteum, produces one of the smallest tree seeds. The great diversity in tree fruits and seeds reflects the many different ways that tree species have evolved to disperse their offspring. For a tree seedling to grow into an adult tree it needs light. If seeds only fell straight to the ground, competition among the concentrated saplings and the shade of the parent would likely prevent it from flourishing. Many seeds such as birch are small and have papery wings to aid dispersal by the wind. Ash trees and maples have larger seeds with blade shaped wings which spiral down to the ground when released. The kapok tree has cottony threads to catch the breeze. The flame tree Delonix regia shoots its seeds through the air when the two sides of its long pods crack apart explosively on drying. The miniature cone-like catkins of alder trees produce seeds that contain small droplets of oil that help disperse the seeds on the surface of water. Mangroves often grow in water and some species have buoyant fruits with seeds that start germinating before they detach from the parent tree. These float on the water and may become lodged on emerging mudbanks and successfully take root.
Other seeds, such as apple pips and plum stones, have fleshy receptacles and smaller fruits like hawthorns have seeds enclosed in edible tissue; animals including mammals and birds eat the fruits and either discard the seeds, or swallow them so they pass through the gut to be deposited in the animal's droppings well away from the parent tree. The germination of some seeds is improved when they are processed in this way. Nuts may be gathered by animals such as squirrels that cache any not immediately consumed. Many of these caches are never revisited; the nut-casing softens with rain and frost, and the surviving seeds germinate in the spring. Pine cones may similarly be hoarded by red squirrels, and grizzly bears may help to disperse the seed by raiding squirrel caches.
The seeds of conifers, the largest group of gymnosperms, are enclosed in a cone and most species have seeds that are light and papery that can be blown considerable distances once free from the cone. Sometimes the seed remains in the cone for years waiting for a trigger event to liberate it. Fire stimulates release and germination of seeds of the jack pine, and also enriches the forest floor with wood ash and removes competing vegetation. Similarly, a number of angiosperms including Acacia cyclops and Acacia mangium have seeds that germinate better after exposure to high temperatures. The single extant species of Ginkgophyta (Ginkgo biloba) has fleshy seeds produced at the ends of short branches on female trees, and Gnetum, a tropical and subtropical group of gymnosperms produce seeds at the tip of a shoot axis.
Evolutionary history
The earliest trees were tree ferns, horsetails and lycophytes, which grew in forests in the Carboniferous period. The first tree may have been Wattieza, fossils of which were found in New York state in 2007 dating back to the Middle Devonian (about 385 million years ago). Prior to this discovery, Archaeopteris was the earliest known tree. Both of these reproduced by spores rather than seeds and are considered to be links between ferns and the gymnosperms which evolved in the Triassic period. The gymnosperms include conifers, cycads, gnetales and ginkgos and these may have appeared as a result of a whole genome duplication event which took place about 319 million years ago. Ginkgophyta was once a widespread diverse group of which the only survivor is the maidenhair tree Ginkgo biloba. This is considered to be a living fossil because it is virtually unchanged from the fossilised specimens found in Triassic deposits.
During the Mesozoic (245 to 66 million years ago) the conifers flourished and became adapted to live in all the major terrestrial habitats. Subsequently, the tree forms of flowering plants evolved during the Cretaceous period. These began to displace the conifers during the Tertiary era (66 to 2 million years ago) when forests covered the globe. When the climate cooled 1.5 million years ago and the first of four glacial periods occurred, the forests retreated as the ice advanced. In the interglacials, trees recolonised the land that had been covered by ice, only to be driven back again in the next glacial period.
Ecology
Trees are an important part of the terrestrial ecosystem, providing essential habitats including many kinds of forest for communities of organisms. Epiphytic plants such as ferns, some mosses, liverworts, orchids and some species of parasitic plants (e.g., mistletoe) hang from branches; these along with arboreal lichens, algae, and fungi provide micro-habitats for themselves and for other organisms, including animals. Leaves, flowers and fruits are seasonally available. On the ground underneath trees there is shade, and often there is undergrowth, leaf litter, and decaying wood that provide other habitat. Trees stabilise the soil, prevent rapid run-off of rain water, help prevent desertification, have a role in climate control and help in the maintenance of biodiversity and ecosystem balance.
Many species of tree support their own specialised invertebrates. In their natural habitats, 284 different species of insect have been found on the English oak (Quercus robur) and 306 species of invertebrate on the Tasmanian oak (Eucalyptus obliqua). Non-native tree species provide a less biodiverse community, for example in the United Kingdom the sycamore (Acer pseudoplatanus), which originates from southern Europe, has few associated invertebrate species, though its bark supports a wide range of lichens, bryophytes and other epiphytes. Trees differ ecologically in the ease with which they can be found by herbivores. Tree apparency varies with a tree's size and semiochemical content, and with the extent to which it is concealed by nonhost neighbours from its insect pests.
In ecosystems such as mangrove swamps, trees play a role in developing the habitat, since the roots of the mangrove trees reduce the speed of flow of tidal currents and trap water-borne sediment, reducing the water depth and creating suitable conditions for further mangrove colonisation. Thus mangrove swamps tend to extend seawards in suitable locations. Mangrove swamps also provide an effective buffer against the more damaging effects of cyclones and tsunamis.
Uses
Food
Trees are the source of many of the world's best known fleshy fruits. Apples, pears, plums, cherries and citrus are all grown commercially in temperate climates and a wide range of edible fruits are found in the tropics. Other commercially important fruit include dates, figs and olives. Palm oil is obtained from the fruits of the oil palm (Elaeis guineensis). The fruits of the cocoa tree (Theobroma cacao) are used to make cocoa and chocolate and the berries of coffee trees, Coffea arabica and Coffea canephora, are processed to extract the coffee beans. In many rural areas of the world, fruit is gathered from forest trees for consumption. Many trees bear edible nuts which can loosely be described as being large, oily kernels found inside a hard shell. These include coconuts (Cocos nucifera), Brazil nuts (Bertholletia excelsa), pecans (Carya illinoinensis), hazel nuts (Corylus), almonds (Prunus dulcis), walnuts (Juglans regia), pistachios (Pistacia vera) and many others. They are high in nutritive value and contain high-quality protein, vitamins and minerals as well as dietary fibre. A variety of nut oils are extracted by pressing for culinary use; some such as walnut, pistachio and hazelnut oils are prized for their distinctive flavours, but they tend to spoil quickly.
In temperate climates there is a sudden movement of sap at the end of the winter as trees prepare to burst into growth. In North America, the sap of the sugar maple (Acer saccharum) is used in the production of maple syrup. About 90% of the sap is water, the remaining 10% being a mixture of various sugars and certain minerals. The sap is harvested by drilling holes in the trunks of the trees and collecting the liquid that flows out of the inserted spigots; the sap is then heated to concentrate the flavour. Similarly in northern Europe the spring rise in the sap of the silver birch (Betula pendula) is tapped and collected, either to be drunk fresh or fermented into an alcoholic drink. In Alaska, the sap of the sweet birch (Betula lenta) is made into a syrup with a sugar content of 67%. Sweet birch sap is more dilute than maple sap; a hundred litres are required to make one litre of birch syrup.
Various parts of trees are used as spices. These include cinnamon, made from the bark of the cinnamon tree (Cinnamomum zeylanicum) and allspice, the dried small fruits of the pimento tree (Pimenta dioica). Nutmeg is a seed found in the fleshy fruit of the nutmeg tree (Myristica fragrans) and cloves are the unopened flower buds of the clove tree (Syzygium aromaticum).
Many trees have flowers rich in nectar which are attractive to bees. The production of forest honey is an important industry in rural areas of the developing world where it is undertaken by small-scale beekeepers using traditional methods. The flowers of the elder (Sambucus) are used to make elderflower cordial and petals of the plum (Prunus spp.) can be candied. Sassafras oil is a flavouring obtained from distilling bark from the roots of the sassafras tree (Sassafras albidum).
The leaves of trees are widely gathered as fodder for livestock and some can be eaten by humans but they tend to be high in tannins which makes them bitter. Leaves of the curry tree (Murraya koenigii) are eaten, those of kaffir lime (Citrus × hystrix) (in Thai food) and Ailanthus (in Korean dishes such as bugak) and those of the European bay tree (Laurus nobilis) and the California bay tree (Umbellularia californica) are used for flavouring food. Camellia sinensis, the source of tea, is a small tree but seldom reaches its full height, being heavily pruned to make picking the leaves easier.
Wood smoke can be used to preserve food. In the hot smoking process the food is exposed to smoke and heat in a controlled environment. The food is ready to eat when the process is complete, having been tenderised and flavoured by the smoke it has absorbed. In the cold process, the temperature is not allowed to rise above . The flavour of the food is enhanced but raw food requires further cooking. If it is to be preserved, meat should be cured before cold smoking.
Fuel
Wood has traditionally been used for fuel, especially in rural areas. In less developed nations it may be the only fuel available and collecting firewood is often a time-consuming task as it becomes necessary to travel further and further afield in the search for fuel. It is often burned inefficiently on an open fire. In more developed countries other fuels are available and burning wood is a choice rather than a necessity. Modern wood-burning stoves are very fuel efficient and new products such as wood pellets are available to burn.
Charcoal can be made by slow pyrolysis of wood by heating it in the absence of air in a kiln. The carefully stacked branches, often oak, are burned with a very limited amount of air. The process of converting them into charcoal takes about fifteen hours. Charcoal is used as a fuel in barbecues and by blacksmiths and has many industrial and other uses.
Timber
Timber, "trees that are grown in order to produce wood" is cut into lumber (sawn wood) for use in construction. Wood has been an important, easily available material for construction since humans started building shelters. Engineered wood products are available which bind the particles, fibres or veneers of wood together with adhesives to form composite materials. Plastics have taken over from wood for some traditional uses.
Wood is used in the construction of buildings, bridges, trackways, piles, poles for power lines, masts for boats, pit props, railway sleepers, fencing, hurdles, shuttering for concrete, pipes, scaffolding and pallets. In housebuilding it is used in joinery, for making joists, roof trusses, roofing shingles, thatching, staircases, doors, window frames, floor boards, parquet flooring, panelling and cladding.
Wood is used to construct carts, farm implements, boats, dugout canoes and in shipbuilding. It is used for making furniture, tool handles, boxes, ladders, musical instruments, bows, weapons, matches, clothes pegs, brooms, shoes, baskets, turnery, carving, toys, pencils, rollers, cogs, wooden screws, barrels, coffins, skittles, veneers, artificial limbs, oars, skis, wooden spoons, sports equipment and wooden balls.
Wood is pulped for paper and used in the manufacture of cardboard and made into engineered wood products for use in construction such as fibreboard, hardboard, chipboard and plywood. The wood of gymnosperms is known as softwood while that of angiosperms is known as hardwood.
Art
Besides inspiring artists down the centuries, trees have been used to create art. Living trees have been used in bonsai and in tree shaping, and both living and dead specimens have been sculpted into sometimes fantastic shapes.
Bonsai
is the practice of growing and shaping small trees, originating in China as penjing and spreading to Japan more than a thousand years ago, there are also similar practices in other cultures like the living miniature landscapes of Vietnam hòn non bộ. The word bonsai is often used in English as an umbrella term for all miniature trees in containers or pots.
The purposes of bonsai are primarily contemplation (for the viewer) and the pleasant exercise of effort and ingenuity (for the grower). Bonsai practice focuses on long-term cultivation and shaping of one or more small trees growing in a container, beginning with a cutting, seedling, or small tree of a species suitable for bonsai development. Bonsai can be created from nearly any perennial woody-stemmed tree or shrub species that produces true branches and can be cultivated to remain small through pot confinement with crown and root pruning. Some species are popular as bonsai material because they have characteristics, such as small leaves or needles, that make them appropriate for the compact visual scope of bonsai and a miniature deciduous forest can even be created using such species as Japanese maple, Japanese zelkova or hornbeam.
Tree shaping
Tree shaping is the practice of changing living trees and other woody plants into man made shapes for art and useful structures. There are a few different methods of shaping a tree. There is a gradual method and there is an instant method. The gradual method slowly guides the growing tip along predetermined pathways over time whereas the instant method bends and weaves saplings long into a shape that becomes more rigid as they thicken up. Most artists use grafting of living trunks, branches, and roots, for art or functional structures and there are plans to grow "living houses" with the branches of trees knitting together to give a solid, weatherproof exterior combined with an interior application of straw and clay to provide a stucco-like inner surface.
Tree shaping has been practised for at least several hundred years, the oldest known examples being the living root bridges built and maintained by the Khasi people of Meghalaya, India using the roots of the rubber tree (Ficus elastica).
Bark
Cork is produced from the thick bark of the cork oak (Quercus suber). It is harvested from the living trees about once every ten years in an environmentally sustainable industry. More than half the world's cork comes from Portugal and is largely used to make stoppers for wine bottles. Other uses include floor tiles, bulletin boards, balls, footwear, cigarette tips, packaging, insulation and joints in woodwind instruments.
The bark of other varieties of oak has traditionally been used in Europe for the tanning of hides though bark from other species of tree has been used elsewhere. The active ingredient, tannin, is extracted and after various preliminary treatments, the skins are immersed in a series of vats containing solutions in increasing concentrations. The tannin causes the hide to become supple, less affected by water and more resistant to bacterial attack.
At least 120 drugs come from plant sources, many of them from the bark of trees. Quinine originates from the cinchona tree (Cinchona) and was for a long time the remedy of choice for the treatment of malaria. Aspirin was synthesised to replace the sodium salicylate derived from the bark of willow trees (Salix) which had unpleasant side effects. The anti-cancer drug Paclitaxel is derived from taxol, a substance found in the bark of the Pacific yew (Taxus brevifolia). Other tree based drugs come from the paw-paw (Carica papaya), the cassia (Cassia spp.), the cocoa tree (Theobroma cacao), the tree of life (Camptotheca acuminata) and the downy birch (Betula pubescens).
The papery bark of the paper birch (Betula papyrifera) tree was used extensively by Native Americans. Wigwams were covered by it and canoes were constructed from it. Other uses included food containers, hunting and fishing equipment, musical instruments, toys and sledges. Nowadays, bark chips, a by-product of the timber industry, are used as a mulch and as a growing medium for epiphytic plants that need a soil-free compost.
Ornamental trees
Trees create a visual impact in the same way as do other landscape features and give a sense of maturity and permanence to park and garden. They are grown for the beauty of their forms, their foliage, flowers, fruit and bark and their siting is of major importance in creating a landscape. They can be grouped informally, often surrounded by plantings of bulbs, laid out in stately avenues or used as specimen trees. As living things, their appearance changes with the season and from year to year.
Trees are often planted in town environments where they are known as street trees or amenity trees. They can provide shade and cooling through evapotranspiration, absorb greenhouse gases and pollutants, intercept rainfall, and reduce the risk of flooding. Scientific studies show that street trees help cities be more sustainable, and improve the physical and mental wellbeing of the citizens. It has been shown that they are beneficial to humans in creating a sense of well-being and reducing stress. Many towns have initiated tree-planting programmes. In London for example, there is an initiative to plant 20,000 new street trees and to have an increase in tree cover of 5% by 2025, equivalent to one tree for every resident.
Other uses
Latex is a sticky defensive secretion that protects plants against herbivores. Many trees produce it when injured but the main source of the latex used to make natural rubber is the Pará rubber tree (Hevea brasiliensis). Originally used to create bouncy balls and for the waterproofing of cloth, natural rubber is now mainly used in tyres for which synthetic materials have proved less durable. The latex exuded by the balatá tree (Manilkara bidentata) is used to make golf balls and is similar to gutta-percha, made from the latex of the "getah perca" tree Palaquium. This is also used as an insulator, particularly of undersea cables, and in dentistry, walking sticks and gun butts. It has now largely been replaced by synthetic materials.
Resin is another plant exudate that may have a defensive purpose. It is a viscous liquid composed mainly of volatile terpenes and is produced mostly by coniferous trees. It is used in varnishes, for making small castings and in ten-pin bowling balls. When heated, the terpenes are driven off and the remaining product is called "rosin" and is used by stringed instrumentalists on their bows. Some resins contain essential oils and are used in incense and aromatherapy. Fossilised resin is known as amber and was mostly formed in the Cretaceous (145 to 66 million years ago) or more recently. The resin that oozed out of trees sometimes trapped insects or spiders and these are still visible in the interior of the amber.
The camphor tree (Cinnamomum camphora) produces an essential oil and the eucalyptus tree (Eucalyptus globulus) is the main source of eucalyptus oil which is used in medicine, as a fragrance and in industry.
Threats
Individual trees
Dead trees pose a safety risk, especially during high winds and severe storms, and removing dead trees involves a financial burden, whereas the presence of healthy trees can clean the air, increase property values, and reduce the temperature of the built environment and thereby reduce building cooling costs. During times of drought, trees can fall into water stress, which may cause a tree to become more susceptible to disease and insect problems, and ultimately may lead to a tree's death. Irrigating trees during dry periods can reduce the risk of water stress and death.
Conservation
About a third of all tree species, some twenty thousand, are included in the IUCN Red List of Threatened Species. Of those, over eight thousand are globally threatened, including at least 1400 which are classed as "critically endangered".
Mythology
Trees have been venerated since time immemorial. To the ancient Celts, certain trees, especially the oak, ash and thorn, held special significance as providing fuel, building materials, ornamental objects and weaponry. Other cultures have similarly revered trees, often linking the lives and fortunes of individuals to them or using them as oracles. In Greek mythology, dryads were believed to be shy nymphs who inhabited trees.
The Oubangui people of west Africa plant a tree when a child is born. As the tree flourishes, so does the child but if the tree fails to thrive, the health of the child is considered at risk. When it flowers it is time for marriage. Gifts are left at the tree periodically and when the individual dies, their spirit is believed to live on in the tree.
Trees have their roots in the ground and their trunk and branches extended towards the sky. This concept is found in many of the world's religions as a tree which links the underworld and the earth and holds up the heavens. In Norse mythology, Yggdrasil is a central cosmic tree whose roots and branches extend to various worlds. Various creatures live on it. In India, Kalpavriksha is a wish-fulfilling tree, one of the nine jewels that emerged from the primitive ocean. Icons are placed beneath it to be worshipped, tree nymphs inhabit the branches and it grants favours to the devout who tie threads round the trunk. Democracy started in North America when the Great Peacemaker formed the Iroquois Confederacy, inspiring the warriors of the original five American nations to bury their weapons under the Tree of Peace, an eastern white pine (Pinus strobus). In the creation story in the Bible, the tree of life and the knowledge of good and evil was planted by God in the Garden of Eden.
Sacred groves exist in China, India, Africa and elsewhere. They are places where the deities live and where all the living things are either sacred or are companions of the gods. Folklore lays down the supernatural penalties that will result if desecration takes place for example by the felling of trees. Because of their protected status, sacred groves may be the only relicts of ancient forest and have a biodiversity much greater than the surrounding area. Some Ancient Indian tree deities, such as Puliyidaivalaiyamman, the Tamil deity of the tamarind tree, or Kadambariyamman, associated with the cadamba tree, were seen as manifestations of a goddess who offers her blessings by giving fruits in abundance.
Superlative trees
Trees have a theoretical maximum height of , but the tallest known specimen on earth is believed to be a coast redwood (Sequoia sempervirens) at Redwood National Park, California. It has been named Hyperion and is tall. In 2006, it was reported to be tall. The tallest known broad-leaved tree is a mountain ash (Eucalyptus regnans) growing in Tasmania with a height of .
The largest tree by volume is believed to be a giant sequoia (Sequoiadendron giganteum) known as the General Sherman Tree in the Sequoia National Park in Tulare County, California. Only the trunk is used in the calculation and the volume is estimated to be .
The oldest living tree with a verified age is also in California. It is a Great Basin bristlecone pine (Pinus longaeva) growing in the White Mountains. It has been dated by drilling a core sample and counting the annual rings. It is estimated to currently be years old.
A little farther south, at Santa Maria del Tule, Oaxaca, Mexico, is the tree with the broadest trunk. It is a Montezuma cypress (Taxodium mucronatum) known as Árbol del Tule and its diameter at breast height is giving it a girth of . The tree's trunk is far from round and the exact dimensions may be misleading as the circumference includes much empty space between the large buttress roots.
| Biology and health sciences | Biology | null |
18955999 | https://en.wikipedia.org/wiki/Desert | Desert | A desert is a landscape where little precipitation occurs and, consequently, living conditions create unique biomes and ecosystems. The lack of vegetation exposes the unprotected surface of the ground to denudation. About one-third of the land surface of the Earth is arid or semi-arid. This includes much of the polar regions, where little precipitation occurs, and which are sometimes called polar deserts or "cold deserts". Deserts can be classified by the amount of precipitation that falls, by the temperature that prevails, by the causes of desertification or by their geographical location.
Deserts are formed by weathering processes as large variations in temperature between day and night put strains on the rocks, which consequently break in pieces. Although rain seldom occurs in deserts, there are occasional downpours that can result in flash floods. Rain falling on hot rocks can cause them to shatter, and the resulting fragments and rubble strewn over the desert floor are further eroded by the wind. This picks up particles of sand and dust, which can remain airborne for extended periods – sometimes causing the formation of sand storms or dust storms. Wind-blown sand grains striking any solid object in their path can abrade the surface. Rocks are smoothed down, and the wind sorts sand into uniform deposits. The grains end up as level sheets of sand or are piled high in billowing sand dunes. Other deserts are flat, stony plains where all the fine material has been blown away and the surface consists of a mosaic of smooth stones, often forming desert pavements, and little further erosion takes place. Other desert features include rock outcrops, exposed bedrock and clays once deposited by flowing water. Temporary lakes may form and salt pans may be left when waters evaporate. There may be underground sources of water, in the form of springs and seepages from aquifers. Where these are found, oases can occur.
Plants and animals living in the desert need special adaptations to survive in the harsh environment. Plants tend to be tough and wiry with small or no leaves, water-resistant cuticles, and often spines to deter herbivory. Some annual plants germinate, bloom and die in the course of a few weeks after rainfall, while other long-lived plants survive for years and have deep root systems able to tap underground moisture. Animals need to keep cool and find enough food and water to survive. Many are nocturnal, and stay in the shade or underground during the heat of the day. They tend to be efficient at conserving water, extracting most of their needs from their food and concentrating their urine. Some animals remain in a state of dormancy for long periods, ready to become active again during the rare rainfall. They then reproduce rapidly while conditions are favorable before returning to dormancy.
People have struggled to live in deserts and the surrounding semi-arid lands for millennia. Nomads have moved their flocks and herds to wherever grazing is available, and oases have provided opportunities for a more settled way of life. The cultivation of semi-arid regions encourages erosion of soil and is one of the causes of increased desertification. Desert farming is possible with the aid of irrigation, and the Imperial Valley in California provides an example of how previously barren land can be made productive by the import of water from an outside source. Many trade routes have been forged across deserts, especially across the Sahara, and traditionally were used by caravans of camels carrying salt, gold, ivory and other goods. Large numbers of slaves were also taken northwards across the Sahara. Some mineral extraction also takes place in deserts, and the uninterrupted sunlight gives potential for the capture of large quantities of solar energy.
Etymology
English desert and its Romance cognates (including Italian and Portuguese deserto, French désert and Spanish desierto) all come from the ecclesiastical Latin dēsertum (originally "an abandoned place"), a participle of dēserere, "to abandon". The correlation between aridity and sparse population is complex and dynamic, varying by culture, era, and technologies; thus the use of the word desert can cause confusion. In English before the 20th century, desert was often used in the sense of "unpopulated area", without specific reference to aridity; but today the word is most often used in its climate-science sense (an area of low precipitation). Phrases such as "desert island" and "Great American Desert", or Shakespeare's "deserts of Bohemia" (The Winter's Tale) in previous centuries did not necessarily imply sand or aridity; their focus was the sparse population.
Major deserts
Deserts occupy about one third of Earth's land surface. Bottomlands may be salt-covered flats. Eolian processes are major factors in shaping desert landscapes. Polar deserts (also seen as "cold deserts") have similar features, except the main form of precipitation is snow rather than rain. Antarctica is the world's largest cold desert (composed of about 98% thick continental ice sheet and 2% barren rock). Some of the barren rock is to be found in the so-called Dry Valleys of Antarctica that almost never get snow, which can have ice-encrusted saline lakes that suggest evaporation far greater than the rare snowfall due to the strong katabatic winds that even evaporate ice.
Deserts, both hot and cold, play a part in moderating Earth's temperature, because they reflect more of the incoming light and their albedo is higher than that of forests or the sea.
Defining characteristics
A desert is a region of land that is very dry because it receives low amounts of precipitation (usually in the form of rain, but it may be snow, mist or fog), often has little coverage by plants, and in which streams dry up unless they are supplied by water from outside the area. Deserts generally receive less than of precipitation each year. The potential evapotranspiration may be large but (in the absence of available water) the actual evapotranspiration may be close to zero. Semi-deserts are regions which receive between and when clad in grass, these are known as steppes. Most deserts on Earth such as the Sahara Desert, Grand Australian Desert and the Great Basin Desert, occur in low altitudes.
Water
One of the driest places on Earth is the Atacama Desert. It is virtually devoid of life because it is blocked from receiving precipitation by the Andes mountains to the east and the Chilean Coast Range to the west. The cold Humboldt Current and the anticyclone of the Pacific are essential to keep the dry climate of the Atacama. The average precipitation in the Chilean region of Antofagasta is just per year. Some weather stations in the Atacama have never received rain. Evidence suggests that the Atacama may not have had any significant rainfall from 1570 to 1971. It is so arid that mountains that reach as high as are completely free of glaciers and, in the southern part from 25°S to 27°S, may have been glacier-free throughout the Quaternary, though permafrost extends down to an altitude of and is continuous above . Nevertheless, there is some plant life in the Atacama, in the form of specialist plants that obtain moisture from dew and the fogs that blow in from the Pacific.
When rain falls in deserts, as it occasionally does, it is often with great violence. The desert surface is evidence of this with dry stream channels known as arroyos or wadis meandering across its surface. These can experience flash floods, becoming raging torrents with surprising rapidity after a storm that may be many kilometers away. Most deserts are in basins with no drainage to the sea but some are crossed by exotic rivers sourced in mountain ranges or other high rainfall areas beyond their borders. The River Nile, the Colorado River and the Yellow River do this, losing much of their water through evaporation as they pass through the desert and raising groundwater levels nearby. There may also be underground sources of water in deserts in the form of springs, aquifers, underground rivers or lakes. Where these lie close to the surface, wells can be dug and oases may form where plant and animal life can flourish. The Nubian Sandstone Aquifer System under the Sahara Desert is the largest known accumulation of fossil water. The Great Man-Made River is a scheme launched by Libya's Muammar Gaddafi to tap this aquifer and supply water to coastal cities. Kharga Oasis in Egypt is long and is the largest oasis in the Libyan Desert. A lake occupied this depression in ancient times and thick deposits of sandy-clay resulted. Wells are dug to extract water from the porous sandstone that lies underneath. Seepages may occur in the walls of canyons and pools may survive in deep shade near the dried up watercourse below.
Lakes may form in basins where there is sufficient precipitation or meltwater from glaciers above. They are usually shallow and saline, and wind blowing over their surface can cause stress, moving the water over nearby low-lying areas. When the lakes dry up, they leave a crust or hardpan behind. This area of deposited clay, silt or sand is known as a playa. The deserts of North America have more than one hundred playas, many of them relics of Lake Bonneville which covered parts of Utah, Nevada and Idaho during the last ice age when the climate was colder and wetter. These include the Great Salt Lake, Utah Lake, Sevier Lake and many dry lake beds. The smooth flat surfaces of playas have been used for attempted vehicle speed records at Black Rock Desert and Bonneville Speedway and the United States Air Force uses Rogers Dry Lake in the Mojave Desert as runways for aircraft and the Space Shuttle.
Classification
Deserts have been defined and classified in a number of ways, generally combining total precipitation, number of days on which this falls, temperature, and humidity, and sometimes additional factors. For example, Phoenix, Arizona, receives less than of precipitation per year, and is immediately recognized as being located in a desert because of its aridity-adapted plants. The North Slope of Alaska's Brooks Range also receives less than of precipitation per year and is often classified as a cold desert. Other regions of the world have cold deserts, including areas of the Himalayas and other high-altitude areas in other parts of the world. Polar deserts cover much of the ice-free areas of the Arctic and Antarctic. A non-technical definition is that deserts are those parts of Earth's surface that have insufficient vegetation cover to support a human population.
Potential evapotranspiration supplements the measurement of precipitation in providing a scientific measurement-based definition of a desert. The water budget of an area can be calculated using the formula P − PE ± S, wherein P is precipitation, PE is potential evapotranspiration rates and S is the amount of surface storage of water. Evapotranspiration is the combination of water loss through atmospheric evaporation and through the life processes of plants. Potential evapotranspiration, then, is the amount of water that could evaporate in any given region. As an example, Tucson, Arizona receives about of rain per year, however about of water could evaporate over the course of a year. In other words, about eight times more water could evaporate from the region than actually falls as rain. Rates of evapotranspiration in cold regions such as Alaska are much lower because of the lack of heat to aid in the evaporation process.
Deserts are sometimes classified as "hot" or "cold", "semiarid" or "coastal". The characteristics of hot deserts include high temperatures in summer; greater evaporation than precipitation, usually exacerbated by high temperatures, strong winds and lack of cloud cover; considerable variation in the occurrence of precipitation, its intensity and distribution; and low humidity. Winter temperatures vary considerably between different deserts and are often related to the location of the desert on the continental landmass and the latitude. Daily variations in temperature can be as great as or more, with heat loss by radiation at night being increased by the clear skies.
Cold deserts, sometimes known as temperate deserts, occur at higher latitudes than hot deserts, and the aridity is caused by the dryness of the air. Some cold deserts are far from the ocean and others are separated by mountain ranges from the sea, and in both cases, there is insufficient moisture in the air to cause much precipitation. The largest of these deserts are found in Central Asia. Others occur on the eastern side of the Rocky Mountains, the eastern side of the southern Andes and in southern Australia. Polar deserts are a particular class of cold desert. The air is very cold and carries little moisture so little precipitation occurs and what does fall, usually snow, is carried along in the often strong wind and may form blizzards, drifts and dunes similar to those caused by dust and sand in other desert regions. In Antarctica, for example, the annual precipitation is about on the central plateau and some ten times that amount on some major peninsulas.
Based on precipitation alone, hyperarid deserts receive less than of rainfall a year; they have no annual seasonal cycle of precipitation and experience twelve-month periods with no rainfall at all. Arid deserts receive between in a year and semiarid deserts between . However, such factors as the temperature, humidity, rate of evaporation and evapotranspiration, and the moisture storage capacity of the ground have a marked effect on the degree of aridity and the plant and animal life that can be sustained. Rain falling in the cold season may be more effective at promoting plant growth, and defining the boundaries of deserts and the semiarid regions that surround them on the grounds of precipitation alone is problematic.
A semi-arid desert or a steppe is a version of the arid desert with much more rainfall, vegetation and higher humidity. These regions feature a semi-arid climate and are less extreme than regular deserts. Like arid deserts, temperatures can vary greatly in semi deserts. They share some characteristics of a true desert and are usually located at the edge of deserts and continental dry areas. They usually receive precipitation from but this can vary due to evapotranspiration and soil nutrition. Semi-deserts can be found in the high elevations of the Tabernas Desert (and some parts of the Spanish Plateu), The Sahel, The Eurasian Steppe, most of Central Asia, the Western US, most of Northern Mexico, portions of South America (especially in Argentina) and the Australian Outback. They usually feature BSh (hot steppe) or BSk (temperate steppe) in the Köppen climate classification.
Coastal deserts are mostly found on the western edges of continental land masses in regions where cold currents approach the land or cold water upwellings rise from the ocean depths. The cool winds crossing this water pick up little moisture and the coastal regions have low temperatures and very low rainfall, the main precipitation being in the form of fog and dew. The range of temperatures on a daily and annual scale is relatively low, being and respectively in the Atacama Desert. Deserts of this type are often long and narrow and bounded to the east by mountain ranges. They occur in Namibia, Chile, southern California and Baja California. Other coastal deserts influenced by cold currents are found in Western Australia, the Arabian Peninsula and Horn of Africa, and the western fringes of the Sahara.
In 1961, Peveril Meigs divided desert regions on Earth into three categories according to the amount of precipitation they received. In this now widely accepted system, extremely arid lands have at least twelve consecutive months without precipitation, arid lands have less than of annual precipitation, and semiarid lands have a mean annual precipitation of between . Both extremely arid and arid lands are considered to be deserts while semiarid lands are generally referred to as steppes when they are grasslands.
Deserts are also classified, according to their geographical location and dominant weather pattern, as trade wind, mid-latitude, rain shadow, coastal, monsoon, or polar deserts. Trade wind deserts occur either side of the horse latitudes at 30° to 35° North and South. These belts are associated with the subtropical anticyclone and the large-scale descent of dry air. The Sahara Desert is of this type. Mid-latitude deserts occur between 30° and 50° North and South. They are mostly in areas remote from the sea where most of the moisture has already precipitated from the prevailing winds. They include the Tengger and Sonoran Deserts. Monsoon deserts are similar. They occur in regions where large temperature differences occur between sea and land. Moist warm air rises over the land, deposits its water content and circulates back to sea. Further inland, areas receive very little precipitation. The Thar Desert near the India/Pakistan border is of this type.
In some parts of the world, deserts are created by a rain shadow effect. Orographic lift occurs as air masses rise to pass over high ground. In the process they cool and lose much of their moisture by precipitation on the windward slope of the mountain range. When they descend on the leeward side, they warm and their capacity to hold moisture increases so an area with relatively little precipitation occurs. The Taklamakan Desert is an example, lying in the rain shadow of the Himalayas and receiving less than precipitation annually. Other areas are arid by virtue of being a very long way from the nearest available sources of moisture.
Montane deserts are arid places with a very high altitude; the most prominent example is found north of the Himalayas, in the Kunlun Mountains and the Tibetan Plateau. Many locations within this category have elevations exceeding and the thermal regime can be hemiboreal. These places owe their profound aridity (the average annual precipitation is often less than 40 mm or 1.5 in) to being very far from the nearest available sources of moisture and are often in the lee of mountain ranges. Montane deserts are normally cold, or may be scorchingly hot by day and very cold by night as is true of the northeastern slopes of Mount Kilimanjaro.
Polar deserts such as McMurdo Dry Valleys remain ice-free because of the dry katabatic winds that flow downhill from the surrounding mountains. Former desert areas presently in non-arid environments, such as the Sandhills in Nebraska, are known as paleodeserts. In the Köppen climate classification system, deserts are classed as BWh (hot desert) or BWk (temperate desert). In the Thornthwaite climate classification system, deserts would be classified as arid megathermal climates.
Polar desert
Polar deserts are a type of cold desert. While they do not lack water, having a persistent cover of snow and ice, this is merely due to marginal evaporation rates and low precipitation.
The McMurdo dry valleys of Antarctica, which lack water (whether rain, ice, or snow) much like a non-polar desert and even have such desert features as hypersaline lakes and intermittent streams that resemble (except for being frozen at their surfaces) hot or cold deserts for extreme aridity and lack of precipitation of any kind. Extreme winds and not seasonal heat desiccate these nearly-lifeless terrains.
Biological desert
The concept of "biological desert" redefines the concept of desert, without the characteristic of aridity, not lacking water, but instead lacking life. Such places can be so-called "ocean deserts", which are mostly at the centers of gyres, but also hypoxic or anoxic waters such as dead zones.
Morphology
Weathering processes
Deserts usually have a large diurnal and seasonal temperature range, with high daytime temperatures falling sharply at night. The diurnal range may be as much as and the rock surface experiences even greater temperature differentials. During the day the sky is usually clear and most of the sun's radiation reaches the ground, but as soon as the sun sets, the desert cools quickly by radiating heat into space. In hot deserts, the temperature during daytime can exceed in summer and plunge below freezing point at night during winter.
Such large temperature variations have a destructive effect on the exposed rocky surfaces. The repeated fluctuations put a strain on exposed rock and the flanks of mountains crack and shatter. Fragmented strata slide down into the valleys where they continue to break into pieces due to the relentless sun by day and chill by night. Successive strata are exposed to further weathering. The relief of the internal pressure that has built up in rocks that have been underground for aeons can cause them to shatter. Exfoliation also occurs when the outer surfaces of rocks split off in flat flakes. This is believed to be caused by the stresses put on the rock by repeated thermal expansions and contractions which induces fracturing parallel to the original surface. Chemical weathering processes probably play a more important role in deserts than was previously thought. The necessary moisture may be present in the form of dew or mist. Ground water may be drawn to the surface by evaporation and the formation of salt crystals may dislodge rock particles as sand or disintegrate rocks by exfoliation. Shallow caves are sometimes formed at the base of cliffs by this means.
As the desert mountains decay, large areas of shattered rock and rubble occur. The process continues and the end products are either dust or sand. Dust is formed from solidified clay or volcanic deposits whereas sand results from the fragmentation of harder granites, limestone and sandstone. There is a certain critical size (about 0.5 mm) below which further temperature-induced weathering of rocks does not occur and this provides a minimum size for sand grains.
As the mountains are eroded, more and more sand is created. At high wind speeds, sand grains are picked up off the surface and blown along, a process known as saltation. The whirling airborne grains act as a sand blasting mechanism which grinds away solid objects in its path as the kinetic energy of the wind is transferred to the ground. The sand eventually ends up deposited in level areas known as sand-fields or sand-seas, or piled up in dunes.
Features
Many people think of deserts as consisting of extensive areas of billowing sand dunes because that is the way they are often depicted on TV and in films, but deserts do not always look like this. Across the world, around 20% of desert is sand, varying from only 2% in North America to 30% in Australia and over 45% in Central Asia. Where sand does occur, it is usually in large quantities in the form of sand sheets or extensive areas of dunes.
A sand sheet is a near-level, firm expanse of partially consolidated particles in a layer that varies from a few centimeters to a few meters thick. The structure of the sheet consists of thin horizontal layers of coarse silt and very fine to medium grain sand, separated by layers of coarse sand and pea-gravel which are a single grain thick. These larger particles anchor the other particles in place and may also be packed together on the surface so as to form a miniature desert pavement. Small ripples form on the sand sheet when the wind exceeds . They form perpendicular to the wind direction and gradually move across the surface as the wind continues to blow. The distance between their crests corresponds to the average length of jumps made by particles during saltation. The ripples are ephemeral and a change in wind direction causes them to reorganise.
Sand dunes are accumulations of windblown sand piled up in mounds or ridges. They form downwind of copious sources of dry, loose sand and occur when topographic and climatic conditions cause airborne particles to settle. As the wind blows, saltation and creep take place on the windward side of the dune and individual grains of sand move uphill. When they reach the crest, they cascade down the far side. The upwind slope typically has a gradient of 10° to 20° while the lee slope is around 32°, the angle at which loose dry sand will slip. As this wind-induced movement of sand grains takes place, the dune moves slowly across the surface of the ground. Dunes are sometimes solitary, but they are more often grouped together in dune fields. When these are extensive, they are known as sand seas or ergs.
The shape of the dune depends on the characteristics of the prevailing wind. Barchan dunes are produced by strong winds blowing across a level surface and are crescent-shaped with the concave side away from the wind. When there are two directions from which winds regularly blow, a series of long, linear dunes known as seif dunes may form. These also occur parallel to a strong wind that blows in one general direction. Transverse dunes run at a right angle to the prevailing wind direction. Star dunes are formed by variable winds, and have several ridges and slip faces radiating from a central point. They tend to grow vertically; they can reach a height of , making them the tallest type of dune. Rounded mounds of sand without a slip face are the rare dome dunes, found on the upwind edges of sand seas.
In deserts where large amounts of limestone mountains surround a closed basin, such as at White Sands National Park in south-central New Mexico, occasional storm runoff transports dissolved limestone and gypsum into a low-lying pan within the basin where the water evaporates, depositing the gypsum and forming crystals known as selenite. The crystals left behind by this process are eroded by the wind and deposited as vast white dune fields that resemble snow-covered landscapes. These types of dune are rare, and only form in closed arid basins that retain the highly soluble gypsum that would otherwise be washed into the sea.
A large part of the surface area of the world's deserts consists of flat, stone-covered plains dominated by wind erosion. In "eolian deflation", the wind continually removes fine-grained material, which becomes wind-blown sand. This exposes coarser-grained material, mainly pebbles with some larger stones or cobbles, leaving a desert pavement, an area of land overlaid by closely packed smooth stones forming a tessellated mosaic. Different theories exist as to how exactly the pavement is formed. It may be that after the sand and dust is blown away by the wind the stones jiggle themselves into place; alternatively, stones previously below ground may in some way work themselves to the surface. Very little further erosion takes place after the formation of a pavement, and the ground becomes stable. Evaporation brings moisture to the surface by capillary action and calcium salts may be precipitated, binding particles together to form a desert conglomerate. In time, bacteria that live on the surface of the stones accumulate a film of minerals and clay particles, forming a shiny brown coating known as desert varnish.
Other non-sandy deserts consist of exposed outcrops of bedrock, dry soils or aridisols, and a variety of landforms affected by flowing water, such as alluvial fans, sinks or playas, temporary or permanent lakes, and oases. A hamada is a type of desert landscape consisting of a high rocky plateau where the sand has been removed by aeolian processes. Other landforms include plains largely covered by gravels and angular boulders, from which the finer particles have been stripped by the wind. These are called "reg" in the western Sahara, "serir" in the eastern Sahara, "gibber plains" in Australia and "saï" in central Asia. The Tassili Plateau in Algeria is a jumble of eroded sandstone outcrops, canyons, blocks, pinnacles, fissures, slabs and ravines. In some places the wind has carved holes or arches, and in others, it has created mushroom-like pillars narrower at the base than the top. On the Colorado Plateau, it is water that has been the prevailing eroding force. Here, rivers, such as the Colorado, have cut their way over the millennia through the high desert floor, creating canyons that are over a mile (6,000 feet or 1,800 meters) deep in places, exposing strata that are over two billion years old.
Dust storms and sandstorms
Sand and dust storms are natural events that occur in arid regions where the land is not protected by a covering of vegetation. Dust storms usually start in desert margins rather than the deserts themselves where the finer materials have already been blown away. As a steady wind begins to blow, fine particles lying on the exposed ground begin to vibrate. At greater wind speeds, some particles are lifted into the air stream. When they land, they strike other particles which may be jerked into the air in their turn, starting a chain reaction. Once ejected, these particles move in one of three possible ways, depending on their size, shape and density; suspension, saltation or creep. Suspension is only possible for particles less than in diameter. In a dust storm, these fine particles are lifted up and wafted aloft to heights of up to . They reduce visibility and can remain in the atmosphere for days on end, conveyed by the trade winds for distances of up to . Denser clouds of dust can be formed in stronger winds, moving across the land with a billowing leading edge. The sunlight can be obliterated and it may become as dark as night at ground level. In a study of a dust storm in China in 2001, it was estimated that 6.5 million tons of dust were involved, covering an area of . The mean particle size was 1.44 μm. A much smaller scale, short-lived phenomenon can occur in calm conditions when hot air near the ground rises quickly through a small pocket of cooler, low-pressure air above forming a whirling column of particles, a dust devil.
Sandstorms occur with much less frequency than dust storms. They are often preceded by severe dust storms and occur when the wind velocity increases to a point where it can lift heavier particles. These grains of sand, up to about in diameter are jerked into the air but soon fall back to earth, ejecting other particles in the process. Their weight prevents them from being airborne for long and most only travel a distance of a few meters (yards). The sand streams along above the surface of the ground like a fluid, often rising to heights of about . In a really severe steady blow, is about as high as the sand stream can rise as the largest sand grains do not become airborne at all. They are transported by creep, being rolled along the desert floor or performing short jumps.
During a sandstorm, the wind-blown sand particles become electrically charged. Such electric fields, which range in size up to 80 kV/m, can produce sparks and cause interference with telecommunications equipment. They are also unpleasant for humans and can cause headaches and nausea. The electric fields are caused by the collision between airborne particles and by the impacts of saltating sand grains landing on the ground. The mechanism is little understood but the particles usually have a negative charge when their diameter is under 250 μm and a positive one when they are over 500 μm.
Ecology and biogeography
Deserts and semi-deserts are home to ecosystems with low or very low biomass and primary productivity in arid or semi-arid climates. They are mostly found in subtropical high-pressure belts and major continental rain shadows. Primary productivity depends on low densities of small photoautotrophs that sustain a sparse trophic network. Plant growth is limited by rainfall, temperature extremes and desiccating winds. Deserts have strong temporal variability in the availability of resources due to the total amount of annual rainfall and the size of individual rainfall events. Resources are often ephemeral or episodic, and this triggers sporadic animal movements and 'pulse and reserve' or 'boom-bust' ecosystem dynamics. Erosion and sedimentation are high due to the sparse vegetation cover and the activities of large mammals and people. Plants and animals in deserts are mostly adapted to extreme and prolonged water deficits, but their reproductive phenology often responds to short episodes of surplus. Competitive interactions are weak.
Flora
Plants face severe challenges in arid environments. Problems they need to solve include how to obtain enough water, how to avoid being eaten and how to reproduce. Photosynthesis is the key to plant growth. It can only take place during the day as energy from the sun is required, but during the day, many deserts become very hot. Opening stomata to allow in the carbon dioxide necessary for the process causes evapotranspiration, and conservation of water is a top priority for desert vegetation. Some plants have resolved this problem by adopting crassulacean acid metabolism, allowing them to open their stomata during the night to allow CO2 to enter, and close them during the day, or by using C4 carbon fixation.
Many desert plants have reduced the size of their leaves or abandoned them altogether. Cacti are present in both North and South America with a post-Gondwana origin. The genus is desert specialist, and in most species, the leaves have been dispensed with and the chlorophyll displaced into the trunks, the cellular structure of which has been modified to allow them to store water. When rain falls, the water is rapidly absorbed by the shallow roots and retained to allow them to survive until the next downpour, which may be months or years away. The giant saguaro cacti of the Sonoran Desert form "forests", providing shade for other plants and nesting places for desert birds. Saguaro grows slowly but may live for up to two hundred years. The surface of the trunk is folded like a concertina, allowing it to expand, and a large specimen can hold eight tons of water after a good downpour.
Other xerophytic plants have developed similar strategies by a process known as convergent evolution. They limit water loss by reducing the size and number of stomata, by having waxy coatings and hairy or tiny leaves. Some are deciduous, shedding their leaves in the driest season, and others curl their leaves up to reduce transpiration. Others, such as aloes, store water in succulent leaves or stems or in fleshy tubers.
Desert plants maximize water uptake by having shallow roots that spread widely, or by developing long taproots that reach down to deep rock strata for ground water. The saltbush in Australia has succulent leaves and secretes salt crystals, enabling it to live in saline areas. In common with cacti, many have developed spines to ward off browsing animals.
Some desert plants produce seed which lies dormant in the soil until sparked into growth by rainfall. With annuals, such plants grow with great rapidity and may flower and set seed within weeks, aiming to complete their development before the last vestige of water dries up. For perennial plants, reproduction is more likely to be successful if the seed germinates in a shaded position, but not so close to the parent plant as to be in competition with it. Some seed will not germinate until it has been blown about on the desert floor to scarify the seed coat. The seed of the mesquite tree, which grows in deserts in the Americas, is hard and fails to sprout even when planted carefully. When it has passed through the gut of a pronghorn it germinates readily, and the little pile of moist dung provides an excellent start to life well away from the parent tree. The stems and leaves of some plants lower the surface velocity of sand-carrying winds and protect the ground from erosion. Even small fungi and microscopic plant organisms found on the soil surface (so-called cryptobiotic soil) can be a vital link in preventing erosion and providing support for other living organisms. Cold deserts often have high concentrations of salt in the soil. Grasses and low shrubs are the dominant vegetation here and the ground may be covered with lichens. Most shrubs have spiny leaves and shed them in the coldest part of the year.
Fauna
Animals adapted to live in deserts are called xerocoles. There is no evidence that body temperature of mammals and birds is adaptive to the different climates, either of great heat or cold. In fact, with a very few exceptions, their basal metabolic rate is determined by body size, irrespective of the climate in which they live. Many desert animals (and plants) show especially clear evolutionary adaptations for water conservation or heat tolerance and so are often studied in comparative physiology, ecophysiology, and evolutionary physiology. One well-studied example is the specializations of mammalian kidneys shown by desert-inhabiting species. Many examples of convergent evolution have been identified in desert organisms, including between cacti and Euphorbia, kangaroo rats and jerboas, Phrynosoma and Moloch lizards.
Deserts present a very challenging environment for animals. Not only do they require food and water but they also need to keep their body temperature at a tolerable level. In many ways, birds are the ablest to do this of the higher animals. They can move to areas of greater food availability as the desert blooms after local rainfall and can fly to faraway waterholes. In hot deserts, gliding birds can remove themselves from the over-heated desert floor by using thermals to soar in the cooler air at great heights. In order to conserve energy, other desert birds run rather than fly. The cream-colored courser flits gracefully across the ground on its long legs, stopping periodically to snatch up insects. Like other desert birds, it is well-camouflaged by its coloring and can merge into the landscape when stationary. The sandgrouse is an expert at this and nests on the open desert floor dozens of kilometers (miles) away from the waterhole it needs to visit daily. Some small diurnal birds are found in very restricted localities where their plumage matches the color of the underlying surface. The desert lark takes frequent dust baths which ensures that it matches its environment.
Water and carbon dioxide are metabolic end products of oxidation of fats, proteins, and carbohydrates. Oxidising a gram of carbohydrate produces 0.60 grams of water; a gram of protein produces 0.41 grams of water; and a gram of fat produces 1.07 grams of water, making it possible for xerocoles to live with little or no access to drinking water. The kangaroo rat for example makes use of this water of metabolism and conserves water both by having a low basal metabolic rate and by remaining underground during the heat of the day, reducing loss of water through its skin and respiratory system when at rest. Herbivorous mammals obtain moisture from the plants they eat. Species such as the addax antelope, dik-dik, Grant's gazelle and oryx are so efficient at doing this that they apparently never need to drink. The camel is a superb example of a mammal adapted to desert life. It minimizes its water loss by producing concentrated urine and dry dung, and is able to lose 40% of its body weight through water loss without dying of dehydration. Carnivores can obtain much of their water needs from the body fluids of their prey. Many other hot desert animals are nocturnal, seeking out shade during the day or dwelling underground in burrows. At depths of more than , these remain at between regardless of the external temperature. Jerboas, desert rats, kangaroo rats and other small rodents emerge from their burrows at night and so do the foxes, coyotes, jackals and snakes that prey on them. Kangaroos keep cool by increasing their respiration rate, panting, sweating and moistening the skin of their forelegs with saliva. Mammals living in cold deserts have developed greater insulation through warmer body fur and insulating layers of fat beneath the skin. The arctic weasel has a metabolic rate that is two or three times as high as would be expected for an animal of its size. Birds have avoided the problem of losing heat through their feet by not attempting to maintain them at the same temperature as the rest of their bodies, a form of adaptive insulation. The emperor penguin has dense plumage, a downy under layer, an air insulation layer next to the skin and various thermoregulatory strategies to maintain its body temperature in one of the harshest environments on Earth.
Being ectotherms, reptiles are unable to live in cold deserts but are well-suited to hot ones. In the heat of the day in the Sahara, the temperature can rise to . Reptiles cannot survive at this temperature and lizards will be prostrated by heat at . They have few adaptations to desert life and are unable to cool themselves by sweating so they shelter during the heat of the day. In the first part of the night, as the ground radiates the heat absorbed during the day, they emerge and search for prey. Lizards and snakes are the most numerous in arid regions and certain snakes have developed a novel method of locomotion that enables them to move sidewards and navigate high sand-dunes. These include the horned viper of Africa and the sidewinder of North America, evolutionarily distinct but with similar behavioural patterns because of convergent evolution. Many desert reptiles are ambush predators and often bury themselves in the sand, waiting for prey to come within range.
Amphibians might seem unlikely desert-dwellers, because of their need to keep their skins moist and their dependence on water for reproductive purposes. In fact, the few species that are found in this habitat have made some remarkable adaptations. Most of them are fossorial, spending the hot dry months aestivating in deep burrows. While there they shed their skins a number of times and retain the remnants around them as a waterproof cocoon to retain moisture. In the Sonoran Desert, Couch's spadefoot toad spends most of the year dormant in its burrow. Heavy rain is the trigger for emergence and the first male to find a suitable pool calls to attract others. Eggs are laid and the tadpoles grow rapidly as they must reach metamorphosis before the water evaporates. As the desert dries out, the adult toads rebury themselves. The juveniles stay on the surface for a while, feeding and growing, but soon dig themselves burrows. Few make it to adulthood. The water holding frog in Australia has a similar life cycle and may aestivate for as long as five years if no rain falls. The Desert rain frog of Namibia is nocturnal and survives because of the damp sea fogs that roll in from the Atlantic.
Invertebrates, particularly arthropods, have successfully made their homes in the desert. Flies, beetles, ants, termites, locusts, millipedes, scorpions and spiders have hard cuticles which are impervious to water and many of them lay their eggs underground and their young develop away from the temperature extremes at the surface. The Saharan silver ant (Cataglyphis bombycina) uses a heat shock protein in a novel way and forages in the open during brief forays in the heat of the day. The long-legged darkling beetle in Namibia stands on its front legs and raises its carapace to catch the morning mist as condensate, funnelling the water into its mouth. Some arthropods make use of the ephemeral pools that form after rain and complete their life cycle in a matter of days. The desert shrimp does this, appearing "miraculously" in new-formed puddles as the dormant eggs hatch. Others, such as brine shrimps, fairy shrimps and tadpole shrimps, are cryptobiotic and can lose up to 92% of their bodyweight, rehydrating as soon as it rains and their temporary pools reappear.
Human relations
Humans have long made use of deserts as places to live, and more recently have started to exploit them for minerals and energy capture. Deserts play a significant role in human culture with an extensive literature. Deserts can only support a limited population of both humans and animals.
History
People have been living in deserts for millennia. Many, such as the Bushmen in the Kalahari, the Aborigines in Australia and various tribes of North American Indians, were originally hunter-gatherers. They developed skills in the manufacture and use of weapons, animal tracking, finding water, foraging for edible plants and using the things they found in their natural environment to supply their everyday needs. Their self-sufficient skills and knowledge were passed down through the generations by word of mouth. Other cultures developed a nomadic way of life as herders of sheep, goats, cattle, camels, yaks, llamas or reindeer. They travelled over large areas with their herds, moving to new pastures as seasonal and erratic rainfall encouraged new plant growth. They took with them their tents made of cloth or skins draped over poles and their diet included milk, blood and sometimes meat.
The desert nomads were also traders. The Sahara is a very large expanse of land stretching from the Atlantic rim to Egypt. Trade routes were developed linking the Sahel in the south with the fertile Mediterranean region to the north and large numbers of camels were used to carry valuable goods across the desert interior. The Tuareg were traders and the transported goods traditionally included slaves, ivory and gold going northwards and salt going southwards. Berbers with knowledge of the region were employed to guide the caravans between the various oases and wells. Several million slaves may have been taken northwards across the Sahara between the 8th and 18th centuries. Traditional means of overland transport declined with the advent of motor vehicles, shipping and air freight, but caravans still travel along routes between Agadez and Bilma and between Timbuktu and Taoudenni carrying salt from the interior to desert-edge communities.
Round the rims of deserts, where more precipitation occurred and conditions were more suitable, some groups took to cultivating crops. This may have happened when drought caused the death of herd animals, forcing herdsmen to turn to cultivation. With few inputs, they were at the mercy of the weather and may have lived at bare subsistence level. The land they cultivated reduced the area available to nomadic herders, causing disputes over land. The semi-arid fringes of the desert have fragile soils which are at risk of erosion when exposed, as happened in the American Dust Bowl in the 1930s. The grasses that held the soil in place were ploughed under, and a series of dry years caused crop failures, while enormous dust storms blew the topsoil away. Half a million Americans were forced to leave their land in this catastrophe.
Similar damage is being done today to the semi-arid areas that rim deserts and about twelve million hectares of land are being turned to desert each year. Desertification is caused by such factors as drought, climatic shifts, tillage for agriculture, overgrazing and deforestation. Vegetation plays a major role in determining the composition of the soil. In many environments, the rate of erosion and run off increases dramatically with reduced vegetation cover.
Natural resource extraction
Deserts contain substantial mineral resources, sometimes over their entire surface, giving them their characteristic colors. For example, the red of many sand deserts comes from laterite minerals. Geological processes in a desert climate can concentrate minerals into valuable deposits. Leaching by ground water can extract ore minerals and redeposit them, according to the water table, in concentrated form. Similarly, evaporation tends to concentrate minerals in desert lakes, creating dry lake beds or playas rich in minerals. Evaporation can concentrate minerals as a variety of evaporite deposits, including gypsum, sodium nitrate, sodium chloride and borates. Evaporites are found in the US's Great Basin Desert, historically exploited by the "20-mule teams" pulling carts of borax from Death Valley to the nearest railway. A desert especially rich in mineral salts is the Atacama Desert, Chile, where sodium nitrate has been mined for explosives and fertilizer since around 1850. Other desert minerals are copper from Chile, Peru, and Iran, and iron and uranium in Australia. Many other metals, salts and commercially valuable types of rock such as pumice are extracted from deserts around the world.
Oil and gas form on the bottom of shallow seas when micro-organisms decompose under anoxic conditions and later become covered with sediment. Many deserts were at one time the sites of shallow seas and others have had underlying hydrocarbon deposits transported to them by the movement of tectonic plates.
Some major oilfields such as Ghawar are found under the sands of Saudi Arabia. Geologists believe that other oil deposits were formed by aeolian processes in ancient deserts as may be the case with some of the major American oil fields.
Farming
Traditional desert farming systems have long been established in North Africa, irrigation being the key to success in an area where water stress is a limiting factor to growth. Techniques that can be used include drip irrigation, the use of organic residues or animal manures as fertilisers and other traditional agricultural management practices. Once fertility has been built up, further crop production preserves the soil from destruction by wind and other forms of erosion. It has been found that plant growth-promoting bacteria play a role in increasing the resistance of plants to stress conditions and these rhizobacterial suspensions could be inoculated into the soil in the vicinity of the plants. A study of these microbes found that desert farming hampers desertification by establishing islands of fertility allowing farmers to achieve increased yields despite the adverse environmental conditions. A field trial in the Sonoran Desert which exposed the roots of different species of tree to rhizobacteria and the nitrogen fixing bacterium Azospirillum brasilense with the aim of restoring degraded lands was only partially successful.
The Judean Desert was farmed in the 7th century BC during the Iron Age to supply food for desert forts. Native Americans in the south western United States became agriculturalists around 600 AD when seeds and technologies became available from Mexico. They used terracing techniques and grew gardens beside seeps, in moist areas at the foot of dunes, near streams providing flood irrigation and in areas irrigated by extensive specially built canals. The Hohokam tribe constructed over of large canals and maintained them for centuries, an impressive feat of engineering. They grew maize, beans, squash and peppers.
A modern example of desert farming is the Imperial Valley in California, which has high temperatures and average rainfall of just per year. The economy is heavily based on agriculture and the land is irrigated through a network of canals and pipelines sourced entirely from the Colorado River via the All-American Canal. The soil is deep and fertile, being part of the river's flood plains, and what would otherwise have been desert has been transformed into one of the most productive farming regions in California. Other water from the river is piped to urban communities but all this has been at the expense of the river, which below the extraction sites no longer has any above-ground flow during most of the year. Another problem of growing crops in this way is the build-up of salinity in the soil caused by the evaporation of river water. The greening of the desert remains an aspiration and was at one time viewed as a future means for increasing food production for the world's growing population. This prospect has proved false as it disregarded the environmental damage caused elsewhere by the diversion of water for desert project irrigation.
Solar energy capture
Deserts are increasingly seen as sources for solar energy, partly due to low amounts of cloud cover. Many solar power plants have been built in the Mojave Desert such as the Solar Energy Generating Systems and Ivanpah Solar Power Facility. Large swaths of this desert are covered in mirrors.
The potential for generating solar energy from the Sahara Desert is huge, the highest found on the globe. Professor David Faiman of Ben-Gurion University has stated that the technology now exists to supply all of the world's electricity needs from 10% of the Sahara Desert. Desertec Industrial Initiative was a consortium seeking $560 billion to invest in North African solar and wind installations over the next forty years to supply electricity to Europe via cable lines running under the Mediterranean Sea. European interest in the Sahara Desert stems from its two aspects: the almost continual daytime sunshine and plenty of unused land. The Sahara receives more sunshine per acre than any part of Europe. The Sahara Desert also has the empty space totalling hundreds of square miles required to house fields of mirrors for solar plants.
The Negev Desert, Israel, and the surrounding area, including the Arava Valley, receive plenty of sunshine and are generally not arable. This has resulted in the construction of many solar plants. David Faiman has proposed that "giant" solar plants in the Negev could supply all of Israel's needs for electricity.
Warfare
The Arabs were probably the first organized force to conduct successful battles in the desert. By knowing back routes and the locations of oases and by utilizing camels, Muslim Arab forces were able to successfully overcome both Roman and Persian forces in the period 600 to 700 AD during the expansion of the Islamic caliphate.
Many centuries later, both world wars saw fighting in the desert. In the First World War, the Ottoman Turks were engaged with the British regular army in a campaign that spanned the Arabian Peninsula. The Turks were defeated by the British, who had the backing of irregular Arab forces that were seeking to revolt against the Turks in the Hejaz, made famous in T.E. Lawrence's book Seven Pillars of Wisdom.
In the Second World War, the Western Desert Campaign began in Italian Libya. Warfare in the desert offered great scope for tacticians to use the large open spaces without the distractions of casualties among civilian populations. Tanks and armoured vehicles were able to travel large distances unimpeded and land mines were laid in large numbers. However, the size and harshness of the terrain meant that all supplies needed to be brought in from great distances. The victors in a battle would advance and their supply chain would necessarily become longer, while the defeated army could retreat, regroup and resupply. For these reasons, the front line moved back and forth through hundreds of kilometers as each side lost and regained momentum. Its most easterly point was at El Alamein in Egypt, where the Allies decisively defeated the Axis forces in 1942.
In culture
The desert is generally thought of as a barren and empty landscape. It has been portrayed by writers, film-makers, philosophers, artists and critics as a place of extremes, a metaphor for anything from death, war or religion to the primitive past or the desolate future.
There is an extensive literature on the subject of deserts. An early historical account is that of Marco Polo (–1324), who travelled through Central Asia to China, crossing a number of deserts in his twenty four year trek. Some accounts give vivid descriptions of desert conditions, though often accounts of journeys across deserts are interwoven with reflection, as is the case in Charles Montagu Doughty's major work, Travels in Arabia Deserta (1888). Antoine de Saint-Exupéry described both his flying and the desert in Wind, Sand and Stars, and Gertrude Bell travelled extensively in the Arabian desert in the early part of the 20th century, becoming an expert on the subject, writing books and advising the British government on dealing with the Arabs. Another woman explorer was Freya Stark, who travelled alone in the Middle East, visiting Turkey, Arabia, Yemen, Syria, Persia and Afghanistan, writing over twenty books on her experiences. The German naturalist Uwe George spent several years living in deserts, recording his experiences and research in his 1976 book, In the Deserts of this Earth.
The American poet Robert Frost expressed his bleak thoughts in his poem, Desert Places (1933), which ends with the stanza "They cannot scare me with their empty spaces / Between stars – on stars where no human race is. / I have it in me so much nearer home / To scare myself with my own desert places."
Saints associated with the desert include Anthony the Great, also known as "Anthony of the Desert". Pope Benedict XVI linked the metaphorical existence of "internal deserts" with physical and social deserts in his homily inaugurating his papacy: "The external deserts in the world are growing, because the internal deserts have become so vast".
Deserts on other planets
Mars is the only other planet in the Solar System besides Earth on which deserts have been identified. Despite its low surface atmospheric pressure (only 1/100 of that of Earth), the patterns of atmospheric circulation on Mars have formed a sea of circumpolar sand more than 5 million km2 (1.9 million sq mi) in the area, larger than most deserts on Earth. The Martian deserts consist of half-moon dunes in flat areas near the permanent polar ice caps in the north. The smaller dune fields occupy the bottom of many of the craters situated in the Martian polar regions. Examination of the surface of rocks by laser beamed from the Mars Exploration Rover have shown a surface film that resembles the desert varnish found on Earth although it might just be surface dust. The surface of Titan, a moon of Saturn, also has a desert-like surface with dune seas.
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2780135 | https://en.wikipedia.org/wiki/Uninhabited%20island | Uninhabited island | An uninhabited island, desert island, or deserted island, is an island, islet or atoll which lacks permanent human population. Uninhabited islands are often depicted in films or stories about shipwrecked people, and are also used as stereotypes for the idea of "paradise". Some uninhabited islands are protected as nature reserves, and some are privately owned. Devon Island in Canada's far north is the largest uninhabited island in the world.
Small coral atolls or islands usually have no source of fresh water, but occasionally a freshwater lens can be reached with a well.
Terminology
Uninhabited islands are sometimes also called "deserted islands" or "desert islands". In the latter, the adjective desert connotes not desert climate conditions, but rather "desolate and sparsely occupied or unoccupied". The word desert has been "formerly applied more widely to any wild, uninhabited region, including forest-land", and it is this archaic meaning that appears in the phrase "desert island".
The term "desert island" is also commonly used figuratively to refer to objects or behavior in conditions of social isolation and limited material means. Behavior on a desert island is a common thought experiment, for example, "desert island morality".
Biodiversity
Desert islands are partly sheltered from humans, making them havens for a number of fragile wildlife species such as sea turtles and ground-nesting seabirds. Many species of seabirds use them as stopovers on their way or especially for nesting, taking advantage of the (supposed) absence of terrestrial predators such as cats or rats.
However, tons of waste from far away countries accumulate on their beaches from the sea, and the absence of surveillance also makes them desirable spots for poachers of protected species.
Selected uninhabited islands
Amatignak Island, southernmost point of Alaska, US
Appat Island, Greenland
ʻAta, the southernmost island of the Kingdom of Tonga
Auckland Islands in the South Pacific, which are part of New Zealand
Astola Island, Pakistan
A majority of the Barra Isles, the Outer Hebrides, Scotland. The most famous of these is Barra Head.
Blasket Islands in County Kerry, Ireland
Ball's Pyramid, a tall volcanic mountain located close to Lord Howe Island in the South Pacific
Binlang Islet, Taiwan
Most of the Canadian Arctic Archipelago
Bouvet Island in the South Atlantic, the world's most remote island
Caquorobert, Guernsey
Clipperton Island, a Pacific island of France
Coral Sea Islands off the northeastern coast of Australia
De Long Islands in the Arctic Ocean, part of Russia
Desertas Islands, Portugal
Devon Island, the largest uninhabited island in the world at 55,247 km2
Gotska Sandön, Sweden
Heard Island and McDonald Islands, Australia
Many islands within the waters of Hong Kong
Ilha da Queimada Grande "Snake Island", Brazil
Isle Royale in Lake Superior
Keros and other small islands off the coast of Greece
Žut and other islands off the coast of Croatia
Kermadec Islands, part of New Zealand
Korzhin Island in Lake Balkhash
Klein Curaçao, Curaçao
Lampione, Sicily
Lítla Dímun, Faroe Islands
Luci Island in Xiuyu District, Putian, Fujian, China
Oeno Island (Pitcairn Islands, British Overseas Territories)
Ogurchinsky Island in the Caspian Sea
Rat Islands, a volcanic island in the Aleutian Islands
Santa Luzia, Cape Verde
Savage Islands, Portugal
Severny Island, Russia
Shag Rocks (South Georgia)
Surtsey, a volcanic island located south of Iceland
Tetepare Island, the largest uninhabited island in the South Pacific
Topo Islet, Azores, Portugal
Most of the United States Minor Outlying Islands, such as Johnston Atoll, Wake Island and Midway Atoll. Palmyra Atoll has no permanent residents but it has private landholdings that are continuously occupied by temporary residents.
Largest uninhabited islands
Most of the largest uninhabited islands are many kilometers/miles inside the Arctic or Antarctic circles, indicating that the reason for their desertedness is the freezing climate.
In literature and popular culture
The first known novels to be set on a desert island were Hayy ibn Yaqdhan written by Ibn Tufail (1105–1185), followed by Theologus Autodidactus written by Ibn al-Nafis (1213–1288). The protagonists in both (Hayy in Hayy ibn Yaqdhan and Kamil in Theologus Autodidactus) are feral children living in seclusion on a deserted island, until they eventually come in contact with castaways from the outside world who are stranded on the island. The story of Theologus Autodidactus, however, extends beyond the deserted island setting when the castaways take Kamil back to civilization with them.
William Shakespeare's 1610–11 play, The Tempest, uses the idea of being stranded on a desert island as a pretext for the action of the play. Prospero and his daughter Miranda are set adrift by Prospero's treacherous brother Antonio, seeking to become Duke of Milan, and Prospero in turn shipwrecks his brother and other men of sin onto the island.
A Latin translation of Ibn Tufail's Hayy ibn Yaqdhan appeared in 1671, prepared by Edward Pococke the Younger, followed by an English translation by Simon Ockley in 1708, as well as German and Dutch translations. In the late 17th century, Hayy ibn Yaqdhan inspired Robert Boyle, an acquaintance of Pococke, to write his own philosophical novel set on a deserted island, The Aspiring Naturalist. Ibn al-Nafis' Theologus Autodidactus was also eventually translated into English in the early 20th century.
Published in 1719, Daniel Defoe's novel Robinson Crusoe, about a castaway on a desert island, has spawned so many imitations in film, television and radio that its name was used to define a genre, Robinsonade. The novel features Man Friday, Crusoe's personal assistant. It is likely that Defoe took inspiration for Crusoe from a Scottish sailor named Alexander Selkirk, who was rescued in 1709 after four years on the otherwise uninhabited Juan Fernández Islands; Defoe usually made use of current events for his plots. It is also likely that he was inspired by the Latin or English translations of Ibn Tufail's Hayy ibn Yaqdhan.
Noel Paul Stookey sang a song about living on a desert island called "On a Desert Island (With You in My Dreams)" on Peter, Paul & Mary's 1965 album: "See What Tomorrow Brings".
Tom Neale was a New Zealander who voluntarily spent 16 years in three sessions in the 1950s and 1960s living alone on the island of Suwarrow in the Northern Cook Islands group. His time there is documented in his autobiography, An Island To Oneself.
In the popular conception, such islands are often located in the Pacific, tropical, uninhabited and usually uncharted. They are remote locales that offer escape and force people marooned or stranded as castaways to become self-sufficient and essentially create a new society. This society can either be utopian, based on an ingenious re-creation of society's comforts (as in Swiss Family Robinson and, in a humorous form, Gilligan's Island) or a regression into savagery (the major theme of both Lord of the Flies and The Beach).
Desert island jokes are also a hugely popular image for gag cartoons, the island being conventionally depicted as just a few yards across with a single palm tree (probably due to the visual constraints of the medium). 17 such cartoons appeared in The New Yorker in 1957 alone.
A special variation of the desert island theme appears in H.G.Wells's The War in the Air. As part of the cataclysmic global war depicted, the bridges linking Goat Island in the middle of the Niagara Falls to the mainland are cut, and with civilization fast breaking down a few survivors stranded on the island cannot expect rescue and must rely on their own resources - embarking on a grim life-and-death struggle.
The top "dream vacation" for heterosexual men surveyed by Psychology Today was "marooned on a tropical island with several members of the opposite sex".
Historical castaways
Essex
In 1820, the crew of the British whaler Essex spent time on uninhabited British Henderson Island. There they gorged on birds, fish, and vegetation and found a small freshwater spring. After one week, they had depleted the island's resources and most of the crew left on three whaleboats, while three of the men decided to remain on the island and survived there for four months until their rescue.
Strathmore
Survivors of the British Strathmore survived for 7 months at a small island of the French Crozet Islands from 1875 to 1876. They survived from eating eggs and flesh of geese, albatrosses and other seabirds. The also ate root vegetables and fish. The survival was the input for among others the book Survival on the Crozet Islands: The Wreck of the Strathmore in 1875.
| Physical sciences | Oceanic and coastal landforms | Earth science |
2781120 | https://en.wikipedia.org/wiki/Hadal%20zone | Hadal zone | The hadal zone, also known as the hadopelagic zone, is the deepest region of the ocean, lying within oceanic trenches. The hadal zone ranges from around below sea level, and exists in long, narrow, topographic V-shaped depressions.
The cumulative area occupied by the 46 individual hadal habitats worldwide is less than 0.25% of the world's seafloor, yet trenches account for over 40% of the ocean's depth range. Most hadal habitat is found in the Pacific Ocean, the deepest of the conventional oceanic divisions.
Terminology and definition
Historically, the hadal zone was not recognized as distinct from the abyssal zone, although the deepest sections were sometimes called "ultra-abyssal". During the early 1950s, the Danish Galathea II and Soviet Vityaz expeditions separately discovered a distinct shift in the life at depths of not recognized by the broad definition of the abyssal zone. The term "hadal" was first proposed in 1956 by Anton Frederik Bruun to describe the parts of the ocean deeper than , leaving abyssal for the parts at . The name refers to Hades, the ancient Greek god of the underworld. About 94% of the hadal zone is found in subduction trenches.
Depths in excess of are generally in ocean trenches, but there are also trenches at shallower depths. These shallower trenches lack the distinct shift in lifeforms and are therefore not hadal. Although the hadal zone has gained widespread recognition and many continue to use the first proposed limit of , it has been observed that represents a gradual transition between the abyssal and hadal zones, leading to the suggestion of placing the limit in the middle, at . Among others, this intermediate limit has been adopted by UNESCO. Similar to other depth ranges, the fauna of the hadal zone can be broadly placed into two groups: the hadobenthic species (compare benthic) living on or at the seabottom/sides of trenches, and the hadopelagic species (compare pelagic) living in the open water.
Ecology
The deepest ocean trenches are considered the least explored and most extreme marine ecosystems. They are characterized by complete lack of sunlight, low temperatures, nutrient scarcity, and extremely high hydrostatic pressures. The major sources of nutrients and carbon are fallout from upper layers, drifts of fine sediment, and landslides. Most organisms are scavengers and detrivores. As of 2020, over 400 species are known from hadal ecosystems, many of which possess physiological adaptations to the extreme environmental conditions. There are high levels of endemism, and noteworthy examples of gigantism in amphipods, mysids, and isopods and dwarfism in nematodes, copepods, and kinorhynchs.
Marine life decreases with depth, both in abundance and biomass, but there is a wide range of metazoan organisms in the hadal zone, mostly benthos, including fish, sea cucumber, bristle worms, bivalves, isopods, sea anemones, amphipods, copepods, decapod crustaceans and gastropods. Most of these trench communities probably originated from the abyssal plains. Although they have evolved adaptations to high pressure and low temperatures such as lower metabolism, intra-cellular protein-stabilising osmolytes, and unsaturated fatty acids in cell membrane phospholipids, there is no consistent relationship between pressure and metabolic rate in these communities. Increased pressure can instead constrain the ontogenic or larval stages of organisms. Pressure increases ten-fold as an organism moves from sea level to a depth of , whilst pressure only doubles as an organism moves from .
Over a geological time scale, trenches can become accessible as previously stenobathic (limited to a narrow depth range) fauna evolve to become eurybathic (adapted to a wider range of depths), such as grenadiers and natantian prawns. Trench communities do, nevertheless, display a contrasting degree of intra-trench endemism and inter-trench similarities at a higher taxonomic level.
Only a relatively small number of fish species are known from the hadal zone, including certain grenadiers, cutthroat eels, pearlfish, cusk-eels, snailfish and eelpouts. Due to the extreme pressure, the theoretical maximum depth for vertebral fish may be about , below which teleosts would be hyperosmotic, assuming trimethylamine N-oxide requirements follow the observed approximate linear relationship with depth. Some invertebrates do occur deeper, such as bigfin squid, certain polynoid worms, myriotrochid sea cucumbers, turrid snails and pardaliscid amphipods in excess of . In addition, giant protists known as Xenophyophora (foraminifera) live at these depths.
Conditions
The only known primary producers in the hadal zone are certain bacteria that are able to metabolize hydrogen and methane released by rock and seawater reactions (serpentinization), or hydrogen sulfide released from cold seeps. Some of these bacteria are symbiotic, for example living inside the mantle of certain thyasirid and vesicomyid bivalves. Otherwise the first link in the hadal food web are heterotroph organisms that feed on marine snow, both fine particles and the occasional carcass.
The hadal zone can reach far below deep; the deepest known extends to . At such depths, the pressure in the hadal zone exceeds . Lack of light and extreme pressure makes this part of the ocean difficult to explore.
Exploration
The exploration of the hadal zone requires the use of instruments that are able to withstand pressures of up to a thousand or more atmospheres. A few haphazard and non-standard tools have been used to collect limited, but valuable, information about the basic biology of a few hadal organisms. Manned and unmanned submersibles, however, can be used to study the depths in greater detail. Unmanned robotic submersibles may be remotely operated (connected to the research vessel by a cable) or autonomous (freely moving). Cameras and manipulators on submersibles allow researchers to observe and take samples of sediment and organisms. Failures of submersibles under the immense pressure at hadal zone depths have occurred. HROV Nereus is thought to have imploded at a depth of 9,990 meters while exploring the Kermadec Trench in 2014.
Notable missions
The first manned exploration to reach Challenger Deep, the deepest known part of the ocean located in the Mariana Trench, was accomplished in 1960 by Jacques Piccard and Don Walsh. They reached a maximum depth of in the bathyscaphe Trieste.
James Cameron also reached the bottom of Mariana Trench in March 2012 using the Deepsea Challenger. The descent of the Deepsea Challenger reached a depth of , slightly less than the deepest dive record set by Piccard and Walsh. Cameron holds the record for the deepest solo dive.
In June 2012, the Chinese manned submersible Jiaolong was able to reach deep in the Mariana Trench, making it the deepest diving manned research submersible. This range surpasses that of the previous record holder, the Japanese-made Shinkai, whose maximum depth is .
Few unmanned submersibles are capable of descending to maximum hadal depths. The deepest diving unmanned submersibles have included the Kaikō (lost at sea in 2003), the ABISMO, the Nereus (lost at sea in 2014), and the Haidou-1.
On July 12, 2022, Dr. Dawn Wright, chief scientist of Esri, completed a scientific expedition to Challenger Deep. Wright served as the mission specialist, with Victor Vescovo piloting the submersible, named the Limiting Factor. This mission made Wright the first African American (of any gender) to reach the Challenger Deep. Additionally, this mission was the first successful side-scan sonar mapping operation at full ocean depth.
| Physical sciences | Oceanography | Earth science |
2782761 | https://en.wikipedia.org/wiki/Southern%20hawker | Southern hawker | The southern hawker or blue hawker (Aeshna cyanea) is a species of hawker dragonfly.
Distribution
The species is one of the most common and most widespread dragonflies in Europe. The range in the Western Palearctic covers a large part of Europe, to Scotland and southern Scandinavia in the north and to Italy (without the southwest) and the northern Balkans in the south); the eastern boundary is formed by the Urals and the western by Ireland. It is also found in northwest Africa (Algeria). In Central Europe the species is very common.
Habitat
These dragonflies mainly inhabit well-vegetated, small ponds and garden ponds, but they wander widely, and they are often seen in gardens and open woodland.
Description
Aeshna cyanea can reach a body length of about , with a wingspan up to . It is a large, brightly coloured dragonfly, with a long body. The thorax is brown, with two ante-humeral wide green longitudinal stripes. On the forehead there is a black spot in the form of the letter T. The wings are hyaline with a dark pterostigma. The leading edge of the wings is dark.
Males have a dark abdomen with bright apple green markings, except for the last segments S8-10 of the abdomen, where the markings are pale blue and joined. In the females, the abdomen is brownish with bright green markings. The eyes of the males have a gradient, fading from green on the bottom to blue or greenish-blue at the top, while in the females they are yellowish green or brownish.
Biology and behavior
Flight period of these insects lasts from June to October, with some specimen visible in May and November. The adults of the southern hawker feed on various insects, caught on the wing. This is an inquisitive species and will approach people.
These dragonflies breed in still or slow-flowing water. The males are often seen patrolling by a ponds edge or river, where they fight away intruders, crashing into rival males and spiralling through the air. The females are quite inconspicuous when they lay their eggs, but they sometimes give away their spot by clattering up from the reeds.
The eggs are laid by jabbing the abdomen into rotting vegetation or wood. The eggs hatch in the spring, after being laid in the previous summer or autumn. The nymphs feed on aquatic insects, small tadpoles, invertebrates and small fish. They emerge as adults in July and August after two to three years’ development.
While looking quite imposing with its relatively large size and voracious appetite for other insects,
the Southern Hawker is completely harmless to humans.
Dragonflies sometimes are observed where the Cercus at the end of their abdomen are mistaken for claws or a stinger, however Dragonflies do not have the ability to sting.
Southern Hawkers, especially the males, are very active fliers and can also be quite inquisitive.
Southern Hawkers can be observed flying hundreds of meters away from their natural habitat flying around open fields, parks, hiking trails, or populated streets, even sometimes flying close to humans, however they generally shy away from getting very close and do not attack humans in any way, unlike wasps or hornets when they feel threatened.
Life cycle
| Biology and health sciences | Odonata | Animals |
2783668 | https://en.wikipedia.org/wiki/Infraspecific%20name | Infraspecific name | In botany, an infraspecific name is the scientific name for any taxon below the rank of species, i.e. an infraspecific taxon or infraspecies. The scientific names of botanical taxa are regulated by the International Code of Nomenclature for algae, fungi, and plants (ICN). As specified by the ICN, the name of an infraspecific taxon is a combination of the name of a species and an infraspecific epithet, separated by a connecting term that denotes the rank of the taxon. An example of an infraspecific name is Astrophytum myriostigma subvar. glabrum, the name of a subvariety of the species Astrophytum myriostigma (bishop's hat cactus). In the previous example, glabrum is the infraspecific epithet.
Names below the rank of species of cultivated kinds of plants and of animals are regulated by different codes of nomenclature and are formed somewhat differently.
Construction of infraspecific names
Article 24 of the ICN describes how infraspecific names are constructed. The order of the three parts of an infraspecific name is:
genus name, specific epithet, connecting term indicating the rank (not part of the name, but required), infraspecific epithet.
It is customary to italicize all three parts of such a name, but not the connecting term. For example:
Acanthocalycium klimpelianum var. macranthum
genus name = Acanthocalycium, specific epithet = klimpelianum, connecting term = var. (short for "varietas" or variety), infraspecific epithet = macranthum
Astrophytum myriostigma subvar. glabrum
genus name = Astrophytum, specific epithet = myriostigma, connecting term = subvar. (short for "subvarietas" or subvariety), infraspecific epithet = glabrum
The recommended abbreviations for ranks below species are:
subspecies - recommended abbreviation: subsp. (but "ssp." is also in use although not recognised by Art 26)
varietas (variety) - recommended abbreviation: var.
subvarietas (subvariety) - recommended abbreviation: subvar.
forma (form) - recommended abbreviation: f.
subforma (subform) - recommended abbreviation: subf.
Although the connecting terms mentioned above are the recommended ones, the ICN allows for other connecting terms in validly published infraspecific taxa. It specifically mentions that Greek letters α, β, γ, etc. can be used in this way in the original document and further ranks may be added without limit. Names that use these connecting terms are now deprecated (though still legal), but they have an importance because they can be basionyms of current species. The commonest cases use "β" and "b"; examples mentioned in the ICN are Cynoglossum cheirifolium β Anchusa (lanata) and Polyporus fomentarius β applanatus whilst other examples (coming from the fungus database Index Fungorum) are Agaricus plexipes b fuliginaria and Peziza capula ß cernua. The ICN allows the possibility that a validly published name could have no defined rank and uses "[unranked]" as the connecting term in such cases.
Abbreviation of infraspecific names
Like specific epithets, infraspecific epithets cannot be used in isolation as names. Thus the name of a particular species of Acanthocalycium is Acanthocalycium klimpelianum, which can be abbreviated to A. klimpelianum where the context makes the genus clear. The species cannot be referred to as just klimpelianum. In the same way, the name of a particular variety of Acanthocalycium klimpelianum is Acanthocalycium klimpelianum var. macranthum, which can be abbreviated to A. k. var. macranthum where the context makes the species clear. The variety cannot be referred to as just macranthum.
Sometimes more than three parts will be given; strictly speaking, this is not a name, but a classification. The ICN gives the example of Saxifraga aizoon var. aizoon subvar. brevifolia f. multicaulis subf. surculosa; the name of the subform would be Saxifraga aizoon subf. surculosa.
Legitimate infraspecific names
For a proposed infraspecific name to be legitimate it must be in accordance with all the rules of the ICN. Only some of the main points are described here.
A key concept in botanical names is that of a type. In many cases the type will be a particular preserved specimen stored in a herbarium, although there are other kinds of type. Like other names, an infraspecific name is attached to a type. Whether a plant should be given a particular infraspecific name can then be decided by comparing it to the type.
There is no requirement for a species to be divided into infraspecific taxa, of whatever rank; in other words, a species does not have to have subspecies, varieties, forms, etc. However, if infraspecific ranks are created, then the name of the type of the species must repeat the specific epithet as its infraspecific epithet. The type acquires this name automatically as soon as any infraspecific rank is created. As an example, consider Poa secunda J.Presl, whose type specimen is in the Wisconsin State Herbarium.
As soon as a subspecies of Poa secunda was created, then the type specimen of P. secunda immediately became the type specimen of Poa secunda subsp. secunda. The name "Poa secunda subsp. secunda" was automatically created (it is an "autonym"). Soreng created the subspecies Poa secunda subsp. juncifolia (whose type specimen is also in the Wisconsin State Herbarium), thereby making the type specimen of P. secunda also the type specimen of Poa secunda subsp. secunda.
If in addition to the subspecies any variety of Poa secunda were to be created, then the type specimen of P. secunda would automatically become the type specimen of Poa secunda var. secunda. The type specimen would then have the classification Poa secunda subsp. secunda var. secunda.
The same epithet can be used again within a species, at whatever level, only if the names with the re-used epithet are attached to the same type. Thus there can be a form called Poa secunda f. juncifolia as well as the subspecies Poa secunda subsp. juncifolia if, and only if, the type specimen of Poa secunda f. juncifolia is the same as the type specimen of Poa secunda subsp. juncifolia (in other words, if there is a single type specimen whose classification is Poa secunda subsp. juncifolia f. juncifolia).
If two infraspecific taxa which have different types are accidentally given the same epithet, then a homonym has been created. The earliest published name is the legitimate one and the other must be changed.
Specifying authors
When indicating authors for infraspecific names, it is possible to show either just the author(s) of the final, infraspecific epithet, or the authors of both the specific and the infraspecific epithets. Examples:
Adenia aculeata subsp. inermis de Wilde
This identifies de Wilde as the author who published this name for the subspecies (i.e. who created the epithet inermis). Note that here it was decided not to indicate authority for the species.
Pinus nigra J.F.Arnold subsp. salzmannii (Dunal) Franco
Here, J.F.Arnold is the author who gave the species, European black pine, its botanical name Pinus nigra; Dunal is the author who was the first to publish the epithet salzmanii for this taxon (as the species Pinus salzmanii); Franco is the author who reduced the taxon to a subspecies of Pinus nigra.
Difference from zoological nomenclature
In zoological nomenclature, names of taxa below species rank are formed somewhat differently, using a trinomen or 'trinomial name'. No connecting term is required as there is only one rank below species, the subspecies.
Difference from prokaryotic nomenclature
The Prokaryotic Code was split from the ICN in 1975. This nomenclature only governs one infraspecific rank, the subspecies, but allows a number of infrasubspecific subdivisions to be used. The authorship is to be specified in the form "Bacillus subtilis subsp. spizizenii Nakamura et al. 1999.", i.e. with only the infraspecific author.
Cultivated plants
The ICN does not regulate the names of cultivated plants, of cultivars, i.e. plants specifically created for use in agriculture or horticulture. Such names are regulated by the International Code of Nomenclature for Cultivated Plants (ICNCP).
Although logically below the rank of species (and hence "infraspecific"), a cultivar name may be attached to any scientific name at the genus level or below. The minimum requirement is to specify a genus name. For example, Achillea 'Cerise Queen' is a cultivar; Pinus nigra 'Arnold Sentinel' is a cultivar of the species P. nigra (which is propagated vegetatively, by cloning).
| Biology and health sciences | Taxonomic rank | Biology |
22719936 | https://en.wikipedia.org/wiki/Baryon%20acoustic%20oscillations | Baryon acoustic oscillations | In cosmology, baryon acoustic oscillations (BAO) are fluctuations in the density of the visible baryonic matter (normal matter) of the universe, caused by acoustic density waves in the primordial plasma of the early universe. In the same way that supernovae provide a "standard candle" for astronomical observations, BAO matter clustering provides a "standard ruler" for length scale in cosmology.
The length of this standard ruler is given by the maximum distance the acoustic waves could travel in the primordial plasma before the plasma cooled to the point where it became neutral atoms (the epoch of recombination), which stopped the expansion of the plasma density waves, "freezing" them into place. The length of this standard ruler (≈490 million light years in today's universe) can be measured by looking at the large scale structure of matter using astronomical surveys. BAO measurements help cosmologists understand more about the nature of dark energy (which causes the accelerating expansion of the universe) by constraining cosmological parameters.
Early universe
The early universe consisted of a hot, dense plasma of electrons and baryons (which include protons and neutrons). Photons (light particles) travelling in this universe were essentially trapped, unable to travel for any considerable distance before interacting with the plasma via Thomson scattering. The average distance which a photon could travel before interacting with the plasma is known as the mean free path of the photon. As the universe expanded, the plasma cooled to below 3000 K—a low enough energy such that the electrons and protons in the plasma could combine to form neutral hydrogen atoms. This recombination happened when the universe was around 379,000 years old, or at a redshift of . At this age, the size of BAO bubbles were 450,000 light-years (0.14 Mpc) in radius (490 million light-years today divided by z = 1089). Photons interact to a much lesser degree with neutral matter, and therefore at recombination the universe became transparent to photons, allowing them to decouple from the matter and free-stream through the universe. The cosmic microwave background (CMB) radiation is light that was scattered just before, and emitted by, recombination, now seen with our telescopes as radio waves all over the sky since it is red-shifted. Therefore, when looking at, for example, Wilkinson Microwave Anisotropy Probe (WMAP) data, one is basically looking back in time to see an image of the universe when it was only 379,000 years old.
WMAP indicates (Figure 1) a smooth, homogeneous universe with density anisotropies of 10 parts per million. However, there are large structures and density fluctuations in the present universe. Galaxies, for instance, are a million times more dense than the universe's mean density. The current belief is that the universe was built in a bottom-up fashion, meaning that the small anisotropies of the early universe acted as gravitational seeds for the structure observed today. Overdense regions attract more matter, whereas underdense regions attract less, and thus these small anisotropies, seen in the CMB, became the large scale structures in the universe today.
Cosmic sound
Imagine an overdense region of the primordial plasma. While this region of overdensity gravitationally attracts matter towards it, the heat of photon-matter interactions creates a large amount of outward pressure. These counteracting forces of gravity and pressure created oscillations, comparable to sound waves created in air by pressure differences.
This overdense region contains dark matter, baryons and photons. The pressure results in spherical sound waves of both baryons and photons moving with a speed slightly over half the speed of light outwards from the overdensity. The dark matter interacts only gravitationally, and so it stays at the center of the sound wave, the origin of the overdensity. Before decoupling, the photons and baryons moved outwards together. After decoupling the photons were no longer interacting with the baryonic matter and they diffused away. That relieved the pressure on the system, leaving behind shells of baryonic matter. Out of all those shells, representing different sound waves wavelengths, the resonant shell corresponds to the first one as it is that shell that travels the same distance for all overdensities before decoupling. This radius is often referred to as the sound horizon.
Without the photo-baryon pressure driving the system outwards, the only remaining force on the baryons was gravitational. Therefore, the baryons and dark matter (left behind at the center of the perturbation) formed a configuration which included overdensities of matter both at the original site of the anisotropy and in the shell at the sound horizon for that anisotropy. Such anisotropies eventually became the ripples in matter density that would form galaxies.
Therefore, one would expect to see a greater number of galaxy pairs separated by the sound horizon distance scale than by other length scales. This particular configuration of matter occurred at each anisotropy in the early universe, and therefore the universe is not composed of one sound ripple, but many overlapping ripples. As an analogy, imagine dropping many pebbles into a pond and watching the resulting wave patterns in the water. It is not possible to observe this preferred separation of galaxies on the sound horizon scale by eye, but one can measure this artifact statistically by looking at the separations of large numbers of galaxies.
Standard ruler
The physics of the propagation of the baryon waves in the early universe is fairly simple; as a result cosmologists can predict the size of the sound horizon at the time of recombination. In addition the CMB provides a measurement of this scale to high accuracy. However, in the time between recombination and present day, the universe has been expanding. This expansion is well supported by observations and is one of the foundations of the Big Bang Model. In the late 1990s, observations of supernovae determined that not only is the universe expanding, it is expanding at an increasing rate. A better understanding of the acceleration of the universe, or dark energy, has become one of the most important questions in cosmology today. In order to understand the nature of the dark energy, it is important to have a variety of ways of measuring the acceleration. BAO can add to the body of knowledge about this acceleration by comparing observations of the sound horizon today (using clustering of galaxies) to that of the sound horizon at the time of recombination (using the CMB). Thus BAO provides a measuring stick with which to better understand the nature of the acceleration, completely independent from the supernova technique.
BAO signal in the Sloan Digital Sky Survey
The Sloan Digital Sky Survey (SDSS) is a major multi-spectral imaging and spectroscopic redshift survey using the dedicated 2.5-metre wide-angle SDSS optical telescope at Apache Point Observatory in New Mexico. The goal of this five-year survey was to take images and spectra of millions of celestial objects. The result of compiling the SDSS data is a three-dimensional map of objects in the nearby universe: the SDSS catalog. The SDSS catalog provides a picture of the distribution of matter in a large enough portion of the universe that one can search for a BAO signal by noting whether there is a statistically significant overabundance of galaxies separated by the predicted sound horizon distance.
The SDSS team looked at a sample of 46,748 luminous red galaxies (LRGs), over 3,816 square-degrees of sky (approximately five billion light years in diameter) and out to a redshift of . They analyzed the clustering of these galaxies by calculating a two-point correlation function on the data. The correlation function (ξ) is a function of comoving galaxy separation distance (s) and describes the probability that one galaxy will be found within a given distance of another.
One would expect a high correlation of galaxies at small separation distances (due to the clumpy nature of galaxy formation) and a low correlation at large separation distances. The BAO signal would show up as a bump in the correlation function at a comoving separation equal to the sound horizon. This signal was detected by the SDSS team in 2005. SDSS confirmed the WMAP results that the sound horizon is ~ in today's universe.
In 2023 astronomers using the SDSS catalog as well as the cosmicflow-4 catalog claimed to have found evidence of an individual BAO bubble with a radius containing some of largest structures known – the Boötes supercluster, the Sloan Great Wall, CfA2 Great Wall, and the Hercules–Corona Borealis Great Wall – which they named Ho'oleilana.
Detection in other galaxy surveys
The 2dFGRS collaboration and the SDSS collaboration reported a detection of the BAO signal in the power spectrum at around the same time in 2005. Both teams are credited and recognized for the discovery by the community as evidenced by the 2014 Shaw Prize in Astronomy which was awarded to both groups. Since then, further detections have been reported in the 6dF Galaxy Survey (6dFGS) in 2011, WiggleZ in 2011 and BOSS in 2012.
Dark energy formalism
BAO constraints on dark energy parameters
The BAO in the radial and transverse directions provide measurements of the Hubble parameter and angular diameter distance, respectively. The angular diameter distance and Hubble parameter can include different functions that explain dark energy behavior. These functions have two parameters w0 and w1 and one can constrain them with a chi-square technique.
General relativity and dark energy
In general relativity, the expansion of the universe is parametrized by a scale factor which is related to redshift:
The Hubble parameter, , in terms of the scale factor is:
where is the time-derivative of the scale factor. The Friedmann equations express the expansion of the universe in terms of Newton's gravitational constant, , the mean gauge pressure, , the Universe's density , the curvature, , and the cosmological constant, :
Observational evidence of the acceleration of the universe implies that (at present time) . Therefore, the following are possible explanations:
The universe is dominated by some field or particle that has negative pressure such that the equation of state:
There is a non-zero cosmological constant, .
The Friedmann equations are incorrect since they contain oversimplifications in order to make the general relativistic field equations easier to compute.
In order to differentiate between these scenarios, precise measurements of the Hubble parameter as a function of redshift are needed.
Measured observables of dark energy
The density parameter, , of various components, , of the universe can be expressed as ratios of the density of to the critical density, :
The Friedman equation can be rewritten in terms of the density parameter. For the current prevailing model of the universe, ΛCDM, this equation is as follows:
where m is matter, r is radiation, k is curvature, Λ is dark energy, and w is the equation of state. Measurements of the CMB from WMAP put tight constraints on many of these parameters; however it is important to confirm and further constrain them using an independent method with different systematics.
The BAO signal is a standard ruler such that the length of the sound horizon can be measured as a function of cosmic time. This measures two cosmological distances: the Hubble parameter, , and the angular diameter distance, , as a function of redshift . By measuring the subtended angle, , of the ruler of length , these parameters are determined as follows:
the redshift interval, , can be measured from the data and thus determining the Hubble parameter as a function of redshift:
Therefore, the BAO technique helps constrain cosmological parameters and provide further insight into the nature of dark energy.
| Physical sciences | Physical cosmology | Astronomy |
28787420 | https://en.wikipedia.org/wiki/Pseudoprime | Pseudoprime | A pseudoprime is a probable prime (an integer that shares a property common to all prime numbers) that is not actually prime. Pseudoprimes are classified according to which property of primes they satisfy.
Some sources use the term pseudoprime to describe all probable primes, both composite numbers and actual primes.
Pseudoprimes are of primary importance in public-key cryptography, which makes use of the difficulty of factoring large numbers into their prime factors. Carl Pomerance estimated in 1988 that it would cost $10 million to factor a number with 144 digits, and $100 billion to factor a 200-digit number (the cost today is dramatically lower but still prohibitively high). But finding two large prime numbers as needed for this use is also expensive, so various probabilistic primality tests are used, some of which in rare cases inappropriately deliver composite numbers instead of primes. On the other hand, deterministic primality tests, such as the AKS primality test, do not give false positives; because of this, there are no pseudoprimes with respect to them.
Fermat pseudoprimes
Fermat's little theorem states that if is prime and is coprime to , then is divisible by . For an integer , if a composite integer divides , then is called a Fermat pseudoprime to base . It follows that if is a Fermat pseudoprime to base , then is coprime to . Some sources use variations of this definition, for example to allow only odd numbers to be pseudoprimes.
An integer that is a Fermat pseudoprime to all values of that are coprime to is called a Carmichael number.
Classes
Catalan pseudoprime
Elliptic pseudoprime
Euler pseudoprime
Euler–Jacobi pseudoprime
Fermat pseudoprime
Frobenius pseudoprime
Lucas pseudoprime
Perrin pseudoprime
Somer–Lucas pseudoprime
Strong pseudoprime
| Mathematics | Prime numbers | null |
533365 | https://en.wikipedia.org/wiki/Rift | Rift | In geology, a rift is a linear zone where the lithosphere is being pulled apart and is an example of extensional tectonics. Typical rift features are a central linear downfaulted depression, called a graben, or more commonly a half-graben with normal faulting and rift-flank uplifts mainly on one side. Where rifts remain above sea level they form a rift valley, which may be filled by water forming a rift lake. The axis of the rift area may contain volcanic rocks, and active volcanism is a part of many, but not all, active rift systems.
Major rifts occur along the central axis of most mid-ocean ridges, where new oceanic crust and lithosphere is created along a divergent boundary between two tectonic plates.
Failed rifts are the result of continental rifting that failed to continue to the point of break-up. Typically the transition from rifting to spreading develops at a triple junction where three converging rifts meet over a hotspot. Two of these evolve to the point of seafloor spreading, while the third ultimately fails, becoming an aulacogen.
Geometry
Most rifts consist of a series of separate segments that together form the linear zone characteristic of rifts. The individual rift segments have a dominantly half-graben geometry, controlled by a single basin-bounding fault. Segment lengths vary between rifts, depending on the elastic thickness of the lithosphere.
Areas of thick colder lithosphere, such as the Baikal Rift, have segment lengths in excess of 80 km, while in areas of warmer thin lithosphere, segment lengths may be less than 30 km. Along the axis of the rift the position, and in some cases the polarity (the dip direction), of the main rift bounding fault changes from segment to segment. Segment boundaries often have a more complex structure and generally cross the rift axis at a high angle. These segment boundary zones accommodate the differences in fault displacement between the segments and are therefore known as accommodation zones.
Accommodation zones take various forms, from a simple relay ramp at the overlap between two major faults of the same polarity, to zones of high structural complexity, particularly where the segments have opposite polarity. Accommodation zones may be located where older crustal structures intersect the rift axis. In the Gulf of Suez rift, the Zaafarana accommodation zone is located where a shear zone in the Arabian-Nubian Shield meets the rift.
Rift flanks or shoulders are elevated areas around rifts. Rift shoulders are typically about 70 km wide. Contrary to what was previously thought, elevated passive continental margins (EPCM) such as the Brazilian Highlands, the Scandinavian Mountains and India's Western Ghats, are not rift shoulders.
Rift development
Rift initiation
The formation of rift basins and strain localization reflects rift maturity. At the onset of rifting, the upper part of the lithosphere starts to extend on a series of initially unconnected normal faults, leading to the development of isolated basins. In subaerial rifts, for example, drainage at the onset of rifting is generally internal, with no element of through drainage.
Mature rift stage
As the rift evolves, some of the individual fault segments grow, eventually becoming linked together to form the larger bounding faults. Subsequent extension becomes concentrated on these faults. The longer faults and wider fault spacing leads to more continuous areas of fault-related subsidence along the rift axis. Significant uplift of the rift shoulders develops at this stage, strongly influencing drainage and sedimentation in the rift basins.
During the climax of lithospheric rifting, as the crust is thinned, the Earth's surface subsides and the Moho becomes correspondingly raised. At the same time, the mantle lithosphere becomes thinned, causing a rise of the top of the asthenosphere. This brings high heat flow from the upwelling asthenosphere into the thinning lithosphere, heating the orogenic lithosphere for dehydration melting, typically causing extreme metamorphism at high thermal gradients of greater than 30 °C. The metamorphic products are high to ultrahigh temperature granulites and their associated migmatite and granites in collisional orogens, with possible emplacement of metamorphic core complexes in continental rift zones but oceanic core complexes in spreading ridges. This leads to a kind of orogeneses in extensional settings, which is referred as to rifting orogeny.
Post-rift subsidence
Once rifting ceases, the mantle beneath the rift cools and this is accompanied by a broad area of post-rift subsidence. The amount of subsidence is directly related to the amount of thinning during the rifting phase calculated as the beta factor (initial crustal thickness divided by final crustal thickness), but is also affected by the degree to which the rift basin is filled at each stage, due to the greater density of sediments in contrast to water. The simple 'McKenzie model' of rifting, which considers the rifting stage to be instantaneous, provides a good first order estimate of the amount of crustal thinning from observations of the amount of post-rift subsidence. This has generally been replaced by the 'flexural cantilever model', which takes into account the geometry of the rift faults and the flexural isostasy of the upper part of the crust.
Multiphase rifting
Some rifts show a complex and prolonged history of rifting, with several distinct phases. The North Sea rift shows evidence of several separate rift phases from the Permian through to the Earliest Cretaceous, a period of over 100 million years.
Rifting to break-up
Rifting may lead to continental breakup and formation of oceanic basins. Successful rifting leads to seafloor spreading along a mid-oceanic ridge and a set of conjugate margins separated by an oceanic basin. Rifting may be active, and controlled by mantle convection. It may also be passive, and driven by far-field tectonic forces that stretch the lithosphere. Margin architecture develops due to spatial and temporal relationships between extensional deformation phases. Margin segmentation eventually leads to the formation of rift domains with variations of the Moho topography, including proximal domain with fault-rotated crustal blocks, necking zone with thinning of crustal basement, distal domain with deep sag basins, ocean-continent transition and oceanic domain.
Deformation and magmatism interact during rift evolution. Magma-rich and magma-poor rifted margins may be formed. Magma-rich margins include major volcanic features. Globally, volcanic margins represent the majority of passive continental margins. Magma-starved rifted margins are affected by large-scale faulting and crustal hyperextension. As a consequence, upper mantle peridotites and gabbros are commonly exposed and serpentinized along extensional detachments at the seafloor.
Magmatism
Many rifts are the sites of at least minor magmatic activity, particularly in the early stages of rifting. Alkali basalts and bimodal volcanism are common products of rift-related magmatism.
Recent studies indicate that post-collisional granites in collisional orogens are the product of rifting magmatism at converged plate margins.
Economic importance
The sedimentary rocks associated with continental rifts host important deposits of both minerals and hydrocarbons.
Mineral deposits
SedEx mineral deposits are found mainly in continental rift settings. They form within post-rift sequences when hydrothermal fluids associated with magmatic activity are expelled at the seabed.
Oil and gas
Continental rifts are the sites of significant oil and gas accumulations, such as the Viking Graben and the Gulf of Suez Rift. Thirty percent of giant oil and gas fields are found within such a setting. In 1999 it was estimated that there were 200 billion barrels of recoverable oil reserves hosted in rifts. Source rocks are often developed within the sediments filling the active rift (syn-rift), forming either in a lacustrine environment or in a restricted marine environment, although not all rifts contain such sequences. Reservoir rocks may be developed in pre-rift, syn-rift and post-rift sequences.
Effective regional seals may be present within the post-rift sequence if mudstones or evaporites are deposited. Just over half of estimated oil reserves are found associated with rifts containing marine syn-rift and post-rift sequences, just under a quarter in rifts with a non-marine syn-rift and post-rift, and an eighth in non-marine syn-rift with a marine post-rift.
Examples
The Asunción Rift in Eastern Paraguay
The Canadian Arctic Rift System in northern North America
The East African Rift
The West and Central African Rift System
The Red Sea Rift
The Gulf of California
The Baikal Rift Zone, the bottom of Lake Baikal is the deepest continental rift on the earth.
The Gulf of Suez Rift
Throughout the Basin and Range Province in North America
The Rio Grande Rift in the southwestern US
The rift zone that contains the Gulf of Corinth in Greece
The Reelfoot Rift, an ancient buried failed rift underlying the New Madrid Seismic Zone in the Mississippi embayment
The Rhine Rift, in southwestern Germany, known as the Upper Rhine valley, part of the European Cenozoic Rift System
The Taupō Volcanic Zone in the North Island of New Zealand
The Oslo Graben in Norway
The Ottawa-Bonnechere Graben in Ontario and Quebec
The Northern Cordilleran Volcanic Province in British Columbia, Yukon and Alaska
The West Antarctic Rift System in Antarctica
The Midcontinent Rift System, a late Precambrian rift in central North America
The Midland Valley in Scotland
The Fundy Basin, a Triassic rift basin in southeastern Canada
The Cambay, Kachchh, and Narmada rifts in northwestern Deccan volcanic province of India
| Physical sciences | Tectonics | Earth science |
533487 | https://en.wikipedia.org/wiki/Energy%20development | Energy development | Energy development is the field of activities focused on obtaining sources of energy from natural resources. These activities include the production of renewable, nuclear, and fossil fuel derived sources of energy, and for the recovery and reuse of energy that would otherwise be wasted. Energy conservation and efficiency measures reduce the demand for energy development, and can have benefits to society with improvements to environmental issues.
Societies use energy for transportation, manufacturing, illumination, heating and air conditioning, and communication, for industrial, commercial, agricultural and domestic purposes. Energy resources may be classified as primary resources, where the resource can be used in substantially its original form, or as secondary resources, where the energy source must be converted into a more conveniently usable form. Non-renewable resources are significantly depleted by human use, whereas renewable resources are produced by ongoing processes that can sustain indefinite human exploitation.
Thousands of people are employed in the energy industry. The conventional industry comprises the petroleum industry, the natural gas industry, the electrical power industry, and the nuclear industry. New energy industries include the renewable energy industry, comprising alternative and sustainable manufacture, distribution, and sale of alternative fuels.
Classification of resources
Energy resources may be classified as primary resources, suitable for end use without conversion to another form, or secondary resources, where the usable form of energy required substantial conversion from a primary source. Examples of primary energy resources are wind power, solar power, wood fuel, fossil fuels such as coal, oil and natural gas, and uranium. Secondary resources are those such as electricity, hydrogen, or other synthetic fuels.
Another important classification is based on the time required to regenerate an energy resource. "Renewable" resources are those that recover their capacity in a time significant by human needs. Examples are hydroelectric power or wind power, when the natural phenomena that are the primary source of energy are ongoing and not depleted by human demands. Non-renewable resources are those that are significantly depleted by human usage and that will not recover their potential significantly during human lifetimes. An example of a non-renewable energy source is coal, which does not form naturally at a rate that would support human use.
Fossil fuels
Fossil fuel (primary non-renewable fossil) sources burn coal or hydrocarbon fuels, which are the remains of the decomposition of plants and animals. There are three main types of fossil fuels: coal, petroleum, and natural gas. Another fossil fuel, liquefied petroleum gas (LPG), is principally derived from the production of natural gas. Heat from burning fossil fuel is used either directly for space heating and process heating, or converted to mechanical energy for vehicles, industrial processes, or electrical power generation. These fossil fuels are part of the carbon cycle and allow solar energy stored in the fuel to be released.
The use of fossil fuels in the 18th and 19th century set the stage for the Industrial Revolution.
Fossil fuels make up the bulk of the world's current primary energy sources. In 2005, 81% of the world's energy needs was met from fossil sources. The technology and infrastructure for the use of fossil fuels already exist. Liquid fuels derived from petroleum deliver much usable energy per unit of weight or volume, which is advantageous when compared with lower energy density sources such as batteries. Fossil fuels are currently economical for decentralized energy use.
Energy dependence on imported fossil fuels creates energy security risks for dependent countries. Oil dependence in particular has led to war, funding of radicals, monopolization, and socio-political instability.
Fossil fuels are non-renewable resources, which will eventually decline in production and become exhausted. While the processes that created fossil fuels are ongoing, fuels are consumed far more quickly than the natural rate of replenishment. Extracting fuels becomes increasingly costly as society consumes the most accessible fuel deposits. Extraction of fossil fuels results in environmental degradation, such as the strip mining and mountaintop removal for coal.
Fuel efficiency is a form of thermal efficiency, meaning the efficiency of a process that converts chemical potential energy contained in a carrier fuel into kinetic energy or work. The fuel economy is the energy efficiency of a particular vehicle, is given as a ratio of distance travelled per unit of fuel consumed. Weight-specific efficiency (efficiency per unit weight) may be stated for freight, and passenger-specific efficiency (vehicle efficiency) per passenger. The inefficient atmospheric combustion (burning) of fossil fuels in vehicles, buildings, and power plants contributes to urban heat islands.
Conventional production of oil peaked, conservatively, between 2007 and 2010. In 2010, it was estimated that an investment of $8 trillion in non-renewable resources would be required to maintain current levels of production for 25 years. In 2010, governments subsidized fossil fuels by an estimated $500 billion a year. Fossil fuels are also a source of greenhouse gas emissions, leading to concerns about global warming if consumption is not reduced.
The combustion of fossil fuels leads to the release of pollution into the atmosphere. The fossil fuels are mainly carbon compounds. During combustion, carbon dioxide is released, and also nitrogen oxides, soot and other fine particulates. The carbon dioxide is the main contributor to recent climate change.
Other emissions from fossil fuel power station include sulphur dioxide, carbon monoxide (CO), hydrocarbons, volatile organic compounds (VOC), mercury, arsenic, lead, cadmium, and other heavy metals including traces of uranium.
A typical coal plant generates billions of kilowatt hours of electrical power per year.
Nuclear
Fission
Nuclear power is the use of nuclear fission to generate useful heat and electricity. Fission of uranium produces nearly all economically significant nuclear power. Radioisotope thermoelectric generators form a very small component of energy generation, mostly in specialized applications such as deep space vehicles.
Nuclear power plants, excluding naval reactors, provided about 5.7% of the world's energy and 13% of the world's electricity in 2012.
In 2013, the IAEA report that there are 437 operational nuclear power reactors, in 31 countries, although not every reactor is producing electricity. In addition, there are approximately 140 naval vessels using nuclear propulsion in operation, powered by some 180 reactors. As of 2013, attaining a net energy gain from sustained nuclear fusion reactions, excluding natural fusion power sources such as the Sun, remains an ongoing area of international physics and engineering research. More than 60 years after the first attempts, commercial fusion power production remains unlikely before 2050.
There is an ongoing debate about nuclear power. Proponents, such as the World Nuclear Association, the IAEA and Environmentalists for Nuclear Energy contend that nuclear power is a safe, sustainable energy source that reduces carbon emissions. Opponents contend that nuclear power poses many threats to people and the environment.
Nuclear power plant accidents include the Chernobyl disaster (1986), Fukushima Daiichi nuclear disaster (2011), and the Three Mile Island accident (1979). There have also been some nuclear submarine accidents. In terms of lives lost per unit of energy generated, analysis has determined that nuclear power has caused less fatalities per unit of energy generated than the other major sources of energy generation. Energy production from coal, petroleum, natural gas and hydropower has caused a greater number of fatalities per unit of energy generated due to air pollution and energy accident effects. However, the economic costs of nuclear power accidents is high, and meltdowns can take decades to clean up. The human costs of evacuations of affected populations and lost livelihoods is also significant.
Comparing Nuclear's latent cancer deaths, such as cancer with other energy sources immediate deaths per unit of energy generated(GWeyr). This study does not include fossil fuel related cancer and other indirect deaths created by the use of fossil fuel consumption in its "severe accident" classification, which would be an accident with more than 5 fatalities.
As of 2012, according to the IAEA, worldwide there were 68 civil nuclear power reactors under construction in 15 countries, approximately 28 of which in the People's Republic of China (PRC), with the most recent nuclear power reactor, as of May 2013, to be connected to the electrical grid, occurring on February 17, 2013, in Hongyanhe Nuclear Power Plant in the PRC. In the United States, two new Generation III reactors are under construction at Vogtle. U.S. nuclear industry officials expect five new reactors to enter service by 2020, all at existing plants. In 2013, four aging, uncompetitive, reactors were permanently closed.
Recent experiments in extraction of uranium use polymer ropes that are coated with a substance that selectively absorbs uranium from seawater. This process could make the considerable volume of uranium dissolved in seawater exploitable for energy production. Since ongoing geologic processes carry uranium to the sea in amounts comparable to the amount that would be extracted by this process, in a sense the sea-borne uranium becomes a sustainable resource.
Nuclear power is a low carbon power generation method of producing electricity, with an analysis of the literature on its total life cycle emission intensity finding that it is similar to renewable sources in a comparison of greenhouse gas (GHG) emissions per unit of energy generated. Since the 1970s, nuclear fuel has displaced about 64 gigatonnes of carbon dioxide equivalent (GtCO2-eq) greenhouse gases, that would have otherwise resulted from the burning of oil, coal or natural gas in fossil-fuel power stations.
Nuclear power phase-out and pull-backs
Japan's 2011 Fukushima Daiichi nuclear accident, which occurred in a reactor design from the 1960s, prompted a rethink of nuclear safety and nuclear energy policy in many countries. Germany decided to close all its reactors by 2022, and Italy has banned nuclear power. Following Fukushima, in 2011 the International Energy Agency halved its estimate of additional nuclear generating capacity to be built by 2035.
Fukushima
Following the 2011 Fukushima Daiichi nuclear disaster – the second worst nuclear incident, that displaced 50,000 households after radioactive material leaked into the air, soil and sea, and with subsequent radiation checks leading to bans on some shipments of vegetables and fish – a global public support survey by Ipsos (2011) for energy sources was published and nuclear fission was found to be the least popular
Fission economics
The economics of new nuclear power plants is a controversial subject, since there are diverging views on this topic, and multibillion-dollar investments ride on the choice of an energy source. Nuclear power plants typically have high capital costs for building the plant, but low direct fuel costs. In recent years there has been a slowdown of electricity demand growth and financing has become more difficult, which affects large projects such as nuclear reactors, with very large upfront costs and long project cycles which carry a large variety of risks. In Eastern Europe, a number of long-established projects are struggling to find finance, notably Belene in Bulgaria and the additional reactors at Cernavoda in Romania, and some potential backers have pulled out. Where cheap gas is available and its future supply relatively secure, this also poses a major problem for nuclear projects.
Analysis of the economics of nuclear power must take into account who bears the risks of future uncertainties. To date all operating nuclear power plants were developed by state-owned or regulated utility monopolies where many of the risks associated with construction costs, operating performance, fuel price, and other factors were borne by consumers rather than suppliers. Many countries have now liberalized the electricity market where these risks, and the risk of cheaper competitors emerging before capital costs are recovered, are borne by plant suppliers and operators rather than consumers, which leads to a significantly different evaluation of the economics of new nuclear power plants.
Costs
Costs are likely to go up for currently operating and new nuclear power plants, due to increased requirements for on-site spent fuel management and elevated design basis threats. While first of their kind designs, such as the EPRs under construction are behind schedule and over-budget, of the seven South Korean APR-1400s presently under construction worldwide, two are in S.Korea at the Hanul Nuclear Power Plant and four are at the largest nuclear station construction project in the world as of 2016, in the United Arab Emirates at the planned Barakah nuclear power plant. The first reactor, Barakah-1 is 85% completed and on schedule for grid-connection during 2017.
Two of the four EPRs under construction (in Finland and France) are significantly behind schedule and substantially over cost.
Renewable sources
Renewable energy is generally defined as energy that comes from resources which are naturally replenished on a human timescale such as sunlight, wind, rain, tides, waves and geothermal heat. Renewable energy replaces conventional fuels in four distinct areas: electricity generation, hot water/space heating, motor fuels, and rural (off-grid) energy services.
Including traditional biomass usage, about 19% of global energy consumption is accounted for by renewable resources. Wind powered energy production is being turned to as a prominent renewable energy source, increasing global wind power capacity by 12% in 2021. While not the case for all countries, 58% of sample countries linked renewable energy consumption to have a positive impact on economic growth. At the national level, at least 30 nations around the world already have renewable energy contributing more than 20% of energy supply. National renewable energy markets are projected to continue to grow strongly in the coming decade and beyond.[76]
Unlike other energy sources, renewable energy sources are not as restricted by geography. Additionally deployment of renewable energy is resulting in economic benefits as well as combating climate change. Rural electrification has been researched on multiple sites and positive effects on commercial spending, appliance use, and general activities requiring electricity as energy. Renewable energy growth in at least 38 countries has been driven by the high electricity usage rates. International support for promoting renewable sources like solar and wind have continued grow.
While many renewable energy projects are large-scale, renewable technologies are also suited to rural and remote areas and developing countries, where energy is often crucial in human development. To ensure human development continues sustainably, governments around the world are beginning to research potential ways to implement renewable sources into their countries and economies. For example, the UK Government’s Department for Energy and Climate Change 2050 Pathways created a mapping technique to educate the public on land competition between energy supply technologies. This tool provides users the ability to understand what the limitations and potential their surrounding land and country has in terms of energy production.
Hydroelectricity
Hydroelectricity is electric power generated by hydropower; the force of falling or flowing water. In 2015 hydropower generated 16.6% of the world's total electricity and 70% of all renewable electricity and was expected to increase about 3.1% each year for the following 25 years.
Hydropower is produced in 150 countries, with the Asia-Pacific region generating 32 percent of global hydropower in 2010. China is the largest hydroelectricity producer, with 721 terawatt-hours of production in 2010, representing around 17 percent of domestic electricity use. There are now three hydroelectricity plants larger than 10 GW: the Three Gorges Dam in China, Itaipu Dam across the Brazil/Paraguay border, and Guri Dam in Venezuela.
The cost of hydroelectricity is relatively low, making it a competitive source of renewable electricity. The average cost of electricity from a hydro plant larger than 10 megawatts is 3 to 5 U.S. cents per kilowatt-hour. Hydro is also a flexible source of electricity since plants can be ramped up and down very quickly to adapt to changing energy demands. However, damming interrupts the flow of rivers and can harm local ecosystems, and building large dams and reservoirs often involves displacing people and wildlife. Once a hydroelectric complex is constructed, the project produces no direct waste, and has a considerably lower output level of the greenhouse gas carbon dioxide than fossil fuel powered energy plants.
Wind
Wind power harnesses the power of the wind to propel the blades of wind turbines. These turbines cause the rotation of magnets, which creates electricity. Wind towers are usually built together on wind farms. There are offshore and onshore wind farms. Global wind power capacity has expanded rapidly to 336 GW in June 2014, and wind energy production was around 4% of total worldwide electricity usage, and growing rapidly.
Wind power is widely used in Europe, Asia, and the United States. Several countries have achieved relatively high levels of wind power penetration, such as 21% of stationary electricity production in Denmark, 18% in Portugal, 16% in Spain, 14% in Ireland, and 9% in Germany in 2010. By 2011, at times over 50% of electricity in Germany and Spain came from wind and solar power. As of 2011, 83 countries around the world are using wind power on a commercial basis.
Many of the world's largest onshore wind farms are located in the United States, China, and India. Most of the world's largest offshore wind farms are located in Denmark, Germany and the United Kingdom. The two largest offshore wind farm are currently the 630 MW London Array and Gwynt y Môr.
Solar
Biofuels
A biofuel is a fuel that contains energy from geologically recent carbon fixation. These fuels are produced from living organisms. Examples of this carbon fixation occur in plants and microalgae. These fuels are made by a biomass conversion (biomass refers to recently living organisms, most often referring to plants or plant-derived materials). This biomass can be converted to convenient energy containing substances in three different ways: thermal conversion, chemical conversion, and biochemical conversion. This biomass conversion can result in fuel in solid, liquid, or gas form. This new biomass can be used for biofuels. Biofuels have increased in popularity because of rising oil prices and the need for energy security.
Bioethanol is an alcohol made by fermentation, mostly from carbohydrates produced in sugar or starch crops such as corn or sugarcane. Cellulosic biomass, derived from non-food sources, such as trees and grasses, is also being developed as a feedstock for ethanol production. Ethanol can be used as a fuel for vehicles in its pure form, but it is usually used as a gasoline additive to increase octane and improve vehicle emissions. Bioethanol is widely used in the USA and in Brazil. Current plant design does not provide for converting the lignin portion of plant raw materials to fuel components by fermentation.
Biodiesel is made from vegetable oils and animal fats. Biodiesel can be used as a fuel for vehicles in its pure form, but it is usually used as a diesel additive to reduce levels of particulates, carbon monoxide, and hydrocarbons from diesel-powered vehicles. Biodiesel is produced from oils or fats using transesterification and is the most common biofuel in Europe. However, research is underway on producing renewable fuels from decarboxylation
In 2010, worldwide biofuel production reached 105 billion liters (28 billion gallons US), up 17% from 2009, and biofuels provided 2.7% of the world's fuels for road transport, a contribution largely made up of ethanol and biodiesel. Global ethanol fuel production reached 86 billion liters (23 billion gallons US) in 2010, with the United States and Brazil as the world's top producers, accounting together for 90% of global production. The world's largest biodiesel producer is the European Union, accounting for 53% of all biodiesel production in 2010. As of 2011, mandates for blending biofuels exist in 31 countries at the national level and in 29 states or provinces. The International Energy Agency has a goal for biofuels to meet more than a quarter of world demand for transportation fuels by 2050 to reduce dependence on petroleum and coal.
Geothermal
Geothermal energy is thermal energy generated and stored in the Earth. Thermal energy is the energy that determines the temperature of matter. The geothermal energy of the Earth's crust originates from the original formation of the planet (20%) and from radioactive decay of minerals (80%). The geothermal gradient, which is the difference in temperature between the core of the planet and its surface, drives a continuous conduction of thermal energy in the form of heat from the core to the surface. The adjective geothermal originates from the Greek roots γη (ge), meaning earth, and θερμος (thermos), meaning hot.
Earth's internal heat is thermal energy generated from radioactive decay and continual heat loss from Earth's formation. Temperatures at the core-mantle boundary may reach over 4000 °C (7,200 °F). The high temperature and pressure in Earth's interior cause some rock to melt and solid mantle to behave plastically, resulting in portions of mantle convecting upward since it is lighter than the surrounding rock. Rock and water is heated in the crust, sometimes up to 370 °C (700 °F).
From hot springs, geothermal energy has been used for bathing since Paleolithic times and for space heating since ancient Roman times, but it is now better known for electricity generation. Worldwide, 11,400 megawatts (MW) of geothermal power is online in 24 countries in 2012. An additional 28 gigawatts of direct geothermal heating capacity is installed for district heating, space heating, spas, industrial processes, desalination and agricultural applications in 2010.
Geothermal power is cost effective, reliable, sustainable, and environmentally friendly, but has historically been limited to areas near tectonic plate boundaries. Recent technological advances have dramatically expanded the range and size of viable resources, especially for applications such as home heating, opening a potential for widespread exploitation. Geothermal wells release greenhouse gases trapped deep within the earth, but these emissions are much lower per energy unit than those of fossil fuels. As a result, geothermal power has the potential to help mitigate global warming if widely deployed in place of fossil fuels.
The Earth's geothermal resources are theoretically more than adequate to supply humanity's energy needs, but only a very small fraction may be profitably exploited. Drilling and exploration for deep resources is very expensive. Forecasts for the future of geothermal power depend on assumptions about technology, energy prices, subsidies, and interest rates. Pilot programs like EWEB's customer opt in Green Power Program show that customers would be willing to pay a little more for a renewable energy source like geothermal. But as a result of government assisted research and industry experience, the cost of generating geothermal power has decreased by 25% over the past two decades. In 2001, geothermal energy cost between two and ten US cents per kWh.
Oceanic
Marine Renewable Energy (MRE) or marine power (also sometimes referred to as ocean energy, ocean power, or marine and hydrokinetic energy) refers to the energy carried by the mechanical energy of ocean waves, currents, and tides, shifts in salinity gradients, and ocean temperature differences. MRE has the potential to become a reliable and renewable energy source because of the cyclical nature of the oceans. The movement of water in the world's oceans creates a vast store of kinetic energy or energy in motion. This energy can be harnessed to generate electricity to power homes, transport, and industries.
The term marine energy encompasses both wave power, i.e. power from surface waves, and tidal power, i.e. obtained from the kinetic energy of large bodies of moving water. Offshore wind power is not a form of marine energy, as wind power is derived from the wind, even if the wind turbines are placed over water. The oceans have a tremendous amount of energy and are close to many if not most concentrated populations. Ocean energy has the potential to provide a substantial amount of new renewable energy around the world.
Marine energy technology is in its first stage of development. To be developed, MRE needs efficient methods of storing, transporting, and capturing ocean power, so it can be used where needed. Over the past year, countries around the world have started implementing market strategies for MRE to commercialize. Canada and China introduced incentives, such as feed-in tariffs (FiTs), which are above-market prices for MRE that allow investors and project developers a stable income. Other financial strategies consist of subsidies, grants, and funding from public-private partnerships (PPPs). China alone approved 100 ocean projects in 2019. Portugal and Spain recognize the potential of MRE in accelerating decarbonization, which is fundamental to meeting the goals of the Paris Agreement. Both countries are focusing on solar and offshore wind auctions to attract private investment, ensure cost-effectiveness, and accelerate MRE growth. Ireland sees MRE as a key component to reduce its carbon footprint. The Offshore Renewable Energy Development Plan (OREDP) supports the exploration and development of the country's significant offshore energy potential. Additionally, Ireland has implemented the Renewable Electricity Support Scheme (RESS) which includes auctions designed to provide financial support for communities, increase technology diversity, and guarantee energy security.
However, while research is increasing, there have been concerns associated with threats to marine mammals, habitats, and potential changes to ocean currents. MRE can be a renewable energy source for coastal communities helping their transition from fossil fuel, but researchers are calling for a better understanding of its environmental impacts. Because ocean-energy areas are often isolated from both fishing and sea traffic, these zones may provide shelter from humans and predators for some marine species. MRE devices can be an ideal home for many fish, crayfish, mollusks, and barnacles; and may also indirectly affect seabirds, and marine mammals because they feed on those species. Similarly, such areas may create an "artificial reef effect" by boosting biodiversity nearby. Noise pollution generated from the technology is limited, also causing fish and mammals living in the area of the installation to return. In the most recent State of Science Report about MRE, the authors claim that there is no evidence for fish, mammals, or seabirds to be injured by either collision, noise pollution, or the electromagnetic field. The uncertainty of its environmental impact comes from the low quantity of MRE devices in the ocean today where data is collected.
100% renewable energy
The incentive to use 100% renewable energy, for electricity, transport, or even total primary energy supply globally, has been motivated by global warming and other ecological as well as economic concerns. Renewable energy use has grown much faster than anyone anticipated. The Intergovernmental Panel on Climate Change has said that there are few fundamental technological limits to integrating a portfolio of renewable energy technologies to meet most of total global energy demand. At the national level, at least 30 nations around the world already have renewable energy contributing more than 20% of energy supply. Also, Stephen W. Pacala and Robert H. Socolow have developed a series of "stabilization wedges" that can allow us to maintain our quality of life while avoiding catastrophic climate change, and "renewable energy sources," in aggregate, constitute the largest number of their "wedges."
Mark Z. Jacobson says producing all new energy with wind power, solar power, and hydropower by 2030 is feasible and existing energy supply arrangements could be replaced by 2050. Barriers to implementing the renewable energy plan are seen to be "primarily social and political, not technological or economic". Jacobson says that energy costs with a wind, solar, water system should be similar to today's energy costs.
Similarly, in the United States, the independent National Research Council has noted that "sufficient domestic renewable resources exist to allow renewable electricity to play a significant role in future electricity generation and thus help confront issues related to climate change, energy security, and the escalation of energy costs ... Renewable energy is an attractive option because renewable resources available in the United States, taken collectively, can supply significantly larger amounts of electricity than the total current or projected domestic demand." .
Critics of the "100% renewable energy" approach include Vaclav Smil and James E. Hansen. Smil and Hansen are concerned about the variable output of solar and wind power, but Amory Lovins argues that the electricity grid can cope, just as it routinely backs up nonworking coal-fired and nuclear plants with working ones.
Google spent $30 million on their "Renewable Energy Cheaper than Coal" project to develop renewable energy and stave off catastrophic climate change. The project was cancelled after concluding that a best-case scenario for rapid advances in renewable energy could only result in emissions 55 percent below the fossil fuel projections for 2050.
Increased energy efficiency
Although increasing the efficiency of energy use is not energy development per se, it may be considered under the topic of energy development since it makes existing energy sources available to do work.
Efficient energy use reduces the amount of energy required to provide products and services. For example, insulating a home allows a building to use less heating and cooling energy to maintain a comfortable temperature. Installing fluorescent lamps or natural skylights reduces the amount of energy required for illumination compared to incandescent light bulbs. Compact fluorescent lights use two-thirds less energy and may last 6 to 10 times longer than incandescent lights. Improvements in energy efficiency are most often achieved by adopting an efficient technology or production process.
Reducing energy use may save consumers money, if the energy savings offsets the cost of an energy efficient technology. Reducing energy use reduces emissions. According to the International Energy Agency, improved energy efficiency in buildings, industrial processes and transportation could reduce the global energy demand in 2050 to around 8% smaller than today, but serving an economy more than twice as big and a population of about 2 billion more people.
Energy efficiency and renewable energy are said to be the twin pillars of sustainable energy policy. In many countries energy efficiency is also seen to have a national security benefit because it can be used to reduce the level of energy imports from foreign countries and may slow down the rate at which domestic energy resources are depleted.
It's been discovered "that for OECD countries, wind, geothermal, hydro and nuclear have the lowest hazard rates among energy sources in production".
Transmission
While new sources of energy are only rarely discovered or made possible by new technology, distribution technology continually evolves. The use of fuel cells in cars, for example, is an anticipated delivery technology. This section presents the various delivery technologies that have been important to historic energy development. They all rely in way on the energy sources listed in the previous section.
Shipping and pipelines
Coal, petroleum and their derivatives are delivered by boat, rail, or road. Petroleum and natural gas may also be delivered by pipeline, and coal via a Slurry pipeline. Fuels such as gasoline and LPG may also be delivered via aircraft. Natural gas pipelines must maintain a certain minimum pressure to function correctly. The higher costs of ethanol transportation and storage are often prohibitive.
Wired energy transfer
Electricity grids are the networks used to transmit and distribute power from production source to end user, when the two may be hundreds of kilometres away. Sources include electrical generation plants such as a nuclear reactor, coal burning power plant, etc. A combination of sub-stations and transmission lines are used to maintain a constant flow of electricity. Grids may suffer from transient blackouts and brownouts, often due to weather damage. During certain extreme space weather events solar wind can interfere with transmissions. Grids also have a predefined carrying capacity or load that cannot safely be exceeded. When power requirements exceed what's available, failures are inevitable. To prevent problems, power is then rationed.
Industrialised countries such as Canada, the US, and Australia are among the highest per capita consumers of electricity in the world, which is possible thanks to a widespread electrical distribution network. The US grid is one of the most advanced, although infrastructure maintenance is becoming a problem. CurrentEnergy provides a realtime overview of the electricity supply and demand for California, Texas, and the Northeast of the US. African countries with small scale electrical grids have a correspondingly low annual per capita usage of electricity. One of the most powerful power grids in the world supplies power to the state of Queensland, Australia.
Wireless energy transfer
Wireless power transfer is a process whereby electrical energy is transmitted from a power source to an electrical load that does not have a built-in power source, without the use of interconnecting wires. Currently available technology is limited to short distances and relatively low power level.
Orbiting solar power collectors would require wireless transmission of power to Earth. The proposed method involves creating a large beam of microwave-frequency radio waves, which would be aimed at a collector antenna site on the Earth. Formidable technical challenges exist to ensure the safety and profitability of such a scheme.
Storage
Energy storage is accomplished by devices or physical media that store energy to perform useful operation at a later time. A device that stores energy is sometimes called an accumulator.
All forms of energy are either potential energy (e.g. Chemical, gravitational, electrical energy, temperature differential, latent heat, etc.) or kinetic energy (e.g. momentum). Some technologies provide only short-term energy storage, and others can be very long-term such as power to gas using hydrogen or methane and the storage of heat or cold between opposing seasons in deep aquifers or bedrock. A wind-up clock stores potential energy (in this case mechanical, in the spring tension), a battery stores readily convertible chemical energy to operate a mobile phone, and a hydroelectric dam stores energy in a reservoir as gravitational potential energy. Ice storage tanks store ice (thermal energy in the form of latent heat) at night to meet peak demand for cooling. Fossil fuels such as coal and gasoline store ancient energy derived from sunlight by organisms that later died, became buried and over time were then converted into these fuels. Even food (which is made by the same process as fossil fuels) is a form of energy stored in chemical form.
History
Since prehistory, when humanity discovered fire to warm up and roast food, through the Middle Ages in which populations built windmills to grind the wheat, until the modern era in which nations can get electricity splitting the atom. Man has sought endlessly for energy sources.
Except nuclear, geothermal and tidal, all other energy sources are from current solar isolation or from fossil remains of plant and animal life that relied upon sunlight. Ultimately, solar energy itself is the result of the Sun's nuclear fusion. Geothermal power from hot, hardened rock above the magma of the Earth's core is the result of the decay of radioactive materials present beneath the Earth's crust, and nuclear fission relies on man-made fission of heavy radioactive elements in the Earth's crust; in both cases these elements were produced in supernova explosions before the formation of the Solar System.
Since the beginning of the Industrial Revolution, the question of the future of energy supplies has been of interest. In 1865, William Stanley Jevons published The Coal Question in which he saw that the reserves of coal were being depleted and that oil was an ineffective replacement. In 1914, U.S. Bureau of Mines stated that the total production was . In 1956, Geophysicist M. King Hubbert deduces that U.S. oil production would peak between 1965 and 1970 and that oil production will peak "within half a century" on the basis of 1956 data. In 1989, predicted peak by Colin Campbell In 2004, OPEC estimated, with substantial investments, it would nearly double oil output by 2025
Sustainability
The environmental movement has emphasized sustainability of energy use and development. Renewable energy is sustainable in its production; the available supply will not be diminished for the foreseeable future - millions or billions of years. "Sustainability" also refers to the ability of the environment to cope with waste products, especially air pollution. Sources which have no direct waste products (such as wind, solar, and hydropower) are brought up on this point. With global demand for energy growing, the need to adopt various energy sources is growing. Energy conservation is an alternative or complementary process to energy development. It reduces the demand for energy by using it efficiently.
Resilience
Some observers contend that idea of "energy independence" is an unrealistic and opaque concept. The alternative offer of "energy resilience" is a goal aligned with economic, security, and energy realities. The notion of resilience in energy was detailed in the 1982 book Brittle Power: Energy Strategy for National Security. The authors argued that simply switching to domestic energy would not be secure inherently because the true weakness is the often interdependent and vulnerable energy infrastructure of a country. Key aspects such as gas lines and the electrical power grid are often centralized and easily susceptible to disruption. They conclude that a "resilient energy supply" is necessary for both national security and the environment. They recommend a focus on energy efficiency and renewable energy that is decentralized.
In 2008, former Intel Corporation Chairman and CEO Andrew Grove looked to energy resilience, arguing that complete independence is unfeasible given the global market for energy. He describes energy resilience as the ability to adjust to interruptions in the supply of energy. To that end, he suggests the U.S. make greater use of electricity. Electricity can be produced from a variety of sources. A diverse energy supply will be less affected by the disruption in supply of any one source. He reasons that another feature of electrification is that electricity is "sticky" – meaning the electricity produced in the U.S. is to stay there because it cannot be transported overseas. According to Grove, a key aspect of advancing electrification and energy resilience will be converting the U.S. automotive fleet from gasoline-powered to electric-powered. This, in turn, will require the modernization and expansion of the electrical power grid. As organizations such as The Reform Institute have pointed out, advancements associated with the developing smart grid would facilitate the ability of the grid to absorb vehicles en masse connecting to it to charge their batteries.
Present and future
Extrapolations from current knowledge to the future offer a choice of energy futures. Predictions parallel the Malthusian catastrophe hypothesis. Numerous are complex models based scenarios as pioneered by Limits to Growth. Modeling approaches offer ways to analyze diverse strategies, and hopefully find a road to rapid and sustainable development of humanity. Short term energy crises are also a concern of energy development. Extrapolations lack plausibility, particularly when they predict a continual increase in oil consumption.
Energy production usually requires an energy investment. Drilling for oil or building a wind power plant requires energy. The fossil fuel resources that are left are often increasingly difficult to extract and convert. They may thus require increasingly higher energy investments. If investment is greater than the value of the energy produced by the resource, it is no longer an effective energy source. These resources are no longer an energy source but may be exploited for value as raw materials. New technology may lower the energy investment required to extract and convert the resources, although ultimately basic physics sets limits that cannot be exceeded.
Between 1950 and 1984, as the Green Revolution transformed agriculture around the globe, world grain production increased by 250%. The energy for the Green Revolution was provided by fossil fuels in the form of fertilizers (natural gas), pesticides (oil), and hydrocarbon fueled irrigation. The peaking of world hydrocarbon production (peak oil) may lead to significant changes, and require sustainable methods of production. One vision of a sustainable energy future involves all human structures on the earth's surface (i.e., buildings, vehicles and roads) doing artificial photosynthesis (using sunlight to split water as a source of hydrogen and absorbing carbon dioxide to make fertilizer) efficiently than plants.
With contemporary space industry's economic activity and the related private spaceflight, with the manufacturing industries, that go into Earth's orbit or beyond, delivering them to those regions will require further energy development. Researchers have contemplated space-based solar power for collecting solar power for use on Earth. Space-based solar power has been in research since the early 1970s. Space-based solar power would require construction of collector structures in space. The advantage over ground-based solar power is higher intensity of light, and no weather to interrupt power collection.
Energy technology
Energy technology is an interdisciplinary engineering science having to do with the efficient, safe, environmentally friendly, and economical extraction, conversion, transportation, storage, and use of energy, targeted towards yielding high efficiency whilst skirting side effects on humans, nature, and the environment.
For people, energy is an overwhelming need, and as a scarce resource, it has been an underlying cause of political conflicts and wars. The gathering and use of energy resources can be harmful to local ecosystems and may have global outcomes.
Energy is also the capacity to do work. We can get energy from food. Energy can be of different forms such as kinetic, potential, mechanical, heat, light etc. Energy is required for individuals and the whole society for lighting, heating, cooking, running, industries, operating transportation and so forth. Basically there are two types of energy depending on the source s they are;
1.Renewable Energy Sources
2.Non-Renewable Energy Sources
Interdisciplinary fields
As an interdisciplinary science Energy technology is linked with many interdisciplinary fields in sundry, overlapping ways.
Physics, for thermodynamics and nuclear physics
Chemistry for fuel, combustion, air pollution, flue gas, battery technology and fuel cells.
Electrical engineering
Engineering, often for fluid energy machines such as combustion engines, turbines, pumps and compressors.
Geography, for geothermal energy and exploration for resources.
Mining, for petrochemical and fossil fuels.
Agriculture and forestry, for sources of renewable energy.
Meteorology for wind and solar energy.
Water and Waterways, for hydropower.
Waste management, for environmental impact.
Transportation, for energy-saving transportation systems.
Environmental studies, for studying the effect of energy use and production on the environment, nature and climate change.
(Lighting Technology), for Interior and Exterior Natural as well as Artificial Lighting Design, Installations, and Energy Savings
(Energy Cost/Benefit Analysis), for Simple Payback and Life Cycle Costing of Energy Efficiency/Conservation Measures Recommended
Electrical engineering
Electric power engineering deals with the production and use of electrical energy, which can entail the study of machines such as generators, electric motors and transformers. Infrastructure involves substations and transformer stations, power lines and electrical cable. Load management and power management over networks have meaningful sway on overall energy efficiency. Electric heating is also widely used and researched.
Thermodynamics
Thermodynamics deals with the fundamental laws of energy conversion and is drawn from theoretical Physics.
Thermal and chemical energy
Thermal and chemical energy are intertwined with chemistry and environmental studies. Combustion has to do with burners and chemical engines of all kinds, grates and incinerators along with their energy efficiency, pollution and operational safety.
Exhaust gas purification technology aims to lessen air pollution through sundry mechanical, thermal and chemical cleaning methods. Emission control technology is a field of process and chemical engineering. Boiler technology deals with the design, construction and operation of steam boilers and turbines (also used in nuclear power generation, see below), drawn from applied mechanics and materials engineering.
Energy conversion has to do with internal combustion engines, turbines, pumps, fans and so on, which are used for transportation, mechanical energy and power generation. High thermal and mechanical loads bring about operational safety worries which are dealt with through many branches of applied engineering science.
Nuclear energy
Nuclear technology deals with nuclear power production from nuclear reactors, along with the processing of nuclear fuel and disposal of radioactive waste, drawing from applied nuclear physics, nuclear chemistry and radiation science.
Nuclear power generation has been politically controversial in many countries for several decades but the electrical energy produced through nuclear fission is of worldwide importance. There are high hopes that fusion technologies will one day replace most fission reactors but this is still a research area of nuclear physics.
Renewable energy
Renewable energy has many branches.
Wind power
Wind turbines convert wind energy into electricity by connecting a spinning rotor to a generator. Wind turbines draw energy from atmospheric currents and are designed using aerodynamics along with knowledge taken from mechanical and electrical engineering. The wind passes across the aerodynamic rotor blades, creating an area of higher pressure and an area of lower pressure on either side of the blade. The forces of lift and drag are formed due to the difference in air pressure. The lift force is stronger than the drag force; therefore the rotor, which is connected to a generator, spins. The energy is then created due to the change from the aerodynamic force to the rotation of the generator.
Being recognized as one of the most efficient renewable energy sources, wind power is becoming more and more relevant and used in the world. Wind power does not use any water in the production of energy making it a good source of energy for areas without much water. Wind energy could also be produced even if the climate changes in line with current predictions, as it relies solely on wind.
Geothermal
Deep within the Earth, is an extreme heat producing layer of molten rock called magma. The very high temperatures from the magma heats nearby groundwater. There are various technologies that have been developed in order to benefit from such heat, such as using different types of power plants (dry, flash or binary), heat pumps, or wells. These processes of harnessing the heat incorporate an infrastructure which has in one form or another a turbine which is spun by either the hot water or the steam produced by it. The spinning turbine, being connected to a generator, produces energy. A more recent innovation involves the use of shallow closed-loop systems that pump heat to and from structures by taking advantage of the constant temperature of soil around 10 feet deep.
Hydropower
Hydropower draws mechanical energy from rivers, ocean waves and tides. Civil engineering is used to study and build dams, tunnels, waterways and manage coastal resources through hydrology and geology. A low speed water turbine spun by flowing water can power an electrical generator to produce electricity.
Bioenergy
Bioenergy deals with the gathering, processing and use of biomasses grown in biological manufacturing, agriculture and forestry from which power plants can draw burning fuel. Ethanol, methanol (both controversial) or hydrogen for fuel cells can be had from these technologies and used to generate electricity.
Enabling technologies
Heat pumps and Thermal energy storage are classes of technologies that can enable the utilization of renewable energy sources that would otherwise be inaccessible due to a temperature that is too low for utilization or a time lag between when the energy is available and when it is needed. While enhancing the temperature of available renewable thermal energy, heat pumps have the additional property of leveraging electrical power (or in some cases mechanical or thermal power) by using it to extract additional energy from a low quality source (such as seawater, lake water, the ground, the air, or waste heat from a process).
Thermal storage technologies allow heat or cold to be stored for periods of time ranging from hours or overnight to interseasonal, and can involve storage of sensible energy (i.e. by changing the temperature of a medium) or latent energy (i.e. through phase changes of a medium, such between water and slush or ice). Short-term thermal storages can be used for peak-shaving in district heating or electrical distribution systems. Kinds of renewable or alternative energy sources that can be enabled include natural energy (e.g. collected via solar-thermal collectors, or dry cooling towers used to collect winter's cold), waste energy (e.g. from HVAC equipment, industrial processes or power plants), or surplus energy (e.g. as seasonally from hydropower projects or intermittently from wind farms). The Drake Landing Solar Community (Alberta, Canada) is illustrative. borehole thermal energy storage allows the community to get 97% of its year-round heat from solar collectors on the garage roofs, which most of the heat collected in summer. Types of storages for sensible energy include insulated tanks, borehole clusters in substrates ranging from gravel to bedrock, deep aquifers, or shallow lined pits that are insulated on top. Some types of storage are capable of storing heat or cold between opposing seasons (particularly if very large), and some storage applications require inclusion of a heat pump. Latent heat is typically stored in ice tanks or what are called phase-change materials (PCMs).
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533867 | https://en.wikipedia.org/wiki/Backup | Backup | In information technology, a backup, or data backup is a copy of computer data taken and stored elsewhere so that it may be used to restore the original after a data loss event. The verb form, referring to the process of doing so, is "back up", whereas the noun and adjective form is "backup". Backups can be used to recover data after its loss from data deletion or corruption, or to recover data from an earlier time. Backups provide a simple form of IT disaster recovery; however not all backup systems are able to reconstitute a computer system or other complex configuration such as a computer cluster, active directory server, or database server.
A backup system contains at least one copy of all data considered worth saving. The data storage requirements can be large. An information repository model may be used to provide structure to this storage. There are different types of data storage devices used for copying backups of data that is already in secondary storage onto archive files. There are also different ways these devices can be arranged to provide geographic dispersion, data security, and portability.
Data is selected, extracted, and manipulated for storage. The process can include methods for dealing with live data, including open files, as well as compression, encryption, and de-duplication. Additional techniques apply to enterprise client-server backup. Backup schemes may include dry runs that validate the reliability of the data being backed up. There are limitations and human factors involved in any backup scheme.
Storage
A backup strategy requires an information repository, "a secondary storage space for data" that aggregates backups of data "sources". The repository could be as simple as a list of all backup media (DVDs, etc.) and the dates produced, or could include a computerized index, catalog, or relational database.
The backup data needs to be stored, requiring a backup rotation scheme, which is a system of backing up data to computer media that limits the number of backups of different dates retained separately, by appropriate re-use of the data storage media by overwriting of backups no longer needed. The scheme determines how and when each piece of removable storage is used for a backup operation and how long it is retained once it has backup data stored on it. The 3-2-1 rule can aid in the backup process. It states that there should be at least 3 copies of the data, stored on 2 different types of storage media, and one copy should be kept offsite, in a remote location (this can include cloud storage). 2 or more different media should be used to eliminate data loss due to similar reasons (for example, optical discs may tolerate being underwater while LTO tapes may not, and SSDs cannot fail due to head crashes or damaged spindle motors since they do not have any moving parts, unlike hard drives). An offsite copy protects against fire, theft of physical media (such as tapes or discs) and natural disasters like floods and earthquakes. Physically protected hard drives are an alternative to an offsite copy, but they have limitations like only being able to resist fire for a limited period of time, so an offsite copy still remains as the ideal choice.
Because there is no perfect storage, many backup experts recommend maintaining a second copy on a local physical device, even if the data is also backed up offsite.
Backup methods
Unstructured
An unstructured repository may simply be a stack of tapes, DVD-Rs or external HDDs with minimal information about what was backed up and when. This method is the easiest to implement, but unlikely to achieve a high level of recoverability as it lacks automation.
Full only/System imaging
A repository using this backup method contains complete source data copies taken at one or more specific points in time. Copying system images, this method is frequently used by computer technicians to record known good configurations. However, imaging is generally more useful as a way of deploying a standard configuration to many systems rather than as a tool for making ongoing backups of diverse systems.
Incremental
An incremental backup stores data changed since a reference point in time. Duplicate copies of unchanged data are not copied. Typically a full backup of all files is made once or at infrequent intervals, serving as the reference point for an incremental repository. Subsequently, a number of incremental backups are made after successive time periods. Restores begin with the last full backup and then apply the incrementals.
Some backup systems can create a from a series of incrementals, thus providing the equivalent of frequently doing a full backup. When done to modify a single archive file, this speeds restores of recent versions of files.
Near-CDP
Continuous Data Protection (CDP) refers to a backup that instantly saves a copy of every change made to the data. This allows restoration of data to any point in time and is the most comprehensive and advanced data protection. Near-CDP backup applications—often marketed as "CDP"—automatically take incremental backups at a specific interval, for example every 15 minutes, one hour, or 24 hours. They can therefore only allow restores to an interval boundary. Near-CDP backup applications use journaling and are typically based on periodic "snapshots", read-only copies of the data frozen at a particular point in time.
Near-CDP (except for Apple Time Machine) intent-logs every change on the host system, often by saving byte or block-level differences rather than file-level differences. This backup method differs from simple disk mirroring in that it enables a roll-back of the log and thus a restoration of old images of data. Intent-logging allows precautions for the consistency of live data, protecting self-consistent files but requiring applications "be quiesced and made ready for backup."
Near-CDP is more practicable for ordinary personal backup applications, as opposed to true CDP, which must be run in conjunction with a virtual machine or equivalent and is therefore generally used in enterprise client-server backups.
Software may create copies of individual files such as written documents, multimedia projects, or user preferences, to prevent failed write events caused by power outages, operating system crashes, or exhausted disk space, from causing data loss. A common implementation is an appended ".bak" extension to the file name.
Reverse incremental
A Reverse incremental backup method stores a recent archive file "mirror" of the source data and a series of differences between the "mirror" in its current state and its previous states. A reverse incremental backup method starts with a non-image full backup. After the full backup is performed, the system periodically synchronizes the full backup with the live copy, while storing the data necessary to reconstruct older versions. This can either be done using hard links—as Apple Time Machine does, or using binary diffs.
Differential
A differential backup saves only the data that has changed since the last full backup. This means a maximum of two backups from the repository are used to restore the data. However, as time from the last full backup (and thus the accumulated changes in data) increases, so does the time to perform the differential backup. Restoring an entire system requires starting from the most recent full backup and then applying just the last differential backup.
A differential backup copies files that have been created or changed since the last full backup, regardless of whether any other differential backups have been made since, whereas an incremental backup copies files that have been created or changed since the most recent backup of any type (full or incremental). Changes in files may be detected through a more recent date/time of last modification file attribute, and/or changes in file size. Other variations of incremental backup include multi-level incrementals and block-level incrementals that compare parts of files instead of just entire files.
Storage media
Regardless of the repository model that is used, the data has to be copied onto an archive file data storage medium. The medium used is also referred to as the type of backup destination.
Magnetic tape
Magnetic tape was for a long time the most commonly used medium for bulk data storage, backup, archiving, and interchange. It was previously a less expensive option, but this is no longer the case for smaller amounts of data. Tape is a sequential access medium, so the rate of continuously writing or reading data can be very fast. While tape media itself has a low cost per space, tape drives are typically dozens of times as expensive as hard disk drives and optical drives.
Many tape formats have been proprietary or specific to certain markets like mainframes or a particular brand of personal computer. By 2014 LTO had become the primary tape technology. The other remaining viable "super" format is the IBM 3592 (also referred to as the TS11xx series). The Oracle StorageTek T10000 was discontinued in 2016.
Hard disk
The use of hard disk storage has increased over time as it has become progressively cheaper. Hard disks are usually easy to use, widely available, and can be accessed quickly. However, hard disk backups are close-tolerance mechanical devices and may be more easily damaged than tapes, especially while being transported. In the mid-2000s, several drive manufacturers began to produce portable drives employing ramp loading and accelerometer technology (sometimes termed a "shock sensor"), and by 2010 the industry average in drop tests for drives with that technology showed drives remaining intact and working after a 36-inch non-operating drop onto industrial carpeting. Some manufacturers also offer 'ruggedized' portable hard drives, which include a shock-absorbing case around the hard disk, and claim a range of higher drop specifications. Over a period of years the stability of hard disk backups is shorter than that of tape backups.
External hard disks can be connected via local interfaces like SCSI, USB, FireWire, or eSATA, or via longer-distance technologies like Ethernet, iSCSI, or Fibre Channel. Some disk-based backup systems, via Virtual Tape Libraries or otherwise, support data deduplication, which can reduce the amount of disk storage capacity consumed by daily and weekly backup data.
Optical storage
Optical storage uses lasers to store and retrieve data. Recordable CDs, DVDs, and Blu-ray Discs are commonly used with personal computers and are generally cheap. The capacities and speeds of these discs have typically been lower than hard disks or tapes. Advances in optical media may shrink that gap in the future.
Potential future data losses caused by gradual media degradation can be predicted by measuring the rate of correctable minor data errors, of which consecutively too many increase the risk of uncorrectable sectors. Support for error scanning varies among optical drive vendors.
Many optical disc formats are WORM type, which makes them useful for archival purposes since the data cannot be changed in any way, including by user error and by malware such as ransomware. Moreover, optical discs are not vulnerable to head crashes, magnetism, imminent water ingress or power surges; and, a fault of the drive typically just halts the spinning.
Optical media is modular; the storage controller is not tied to media itself like with hard drives or flash storage (→flash memory controller), allowing it to be removed and accessed through a different drive. However, recordable media may degrade earlier under long-term exposure to light.
Some optical storage systems allow for cataloged data backups without human contact with the discs, allowing for longer data integrity. A French study in 2008 indicated that the lifespan of typically-sold CD-Rs was 2–10 years, but one manufacturer later estimated the longevity of its CD-Rs with a gold-sputtered layer to be as high as 100 years. Sony's proprietary Optical Disc Archive can in 2016 reach a read rate of 250 MB/s.
Solid-state drive
Solid-state drives (SSDs) use integrated circuit assemblies to store data. Flash memory, thumb drives, USB flash drives, CompactFlash, SmartMedia, Memory Sticks, and Secure Digital card devices are relatively expensive for their low capacity, but convenient for backing up relatively low data volumes. A solid-state drive does not contain any movable parts, making it less susceptible to physical damage, and can have huge throughput of around 500 Mbit/s up to 6 Gbit/s. Available SSDs have become more capacious and cheaper. Flash memory backups are stable for fewer years than hard disk backups.
Remote backup service
Remote backup services or cloud backups involve service providers storing data offsite. This has been used to protect against events such as fires, floods, or earthquakes which could destroy locally stored backups. Cloud-based backup (through services like or similar to Google Drive, and Microsoft OneDrive) provides a layer of data protection. However, the users must trust the provider to maintain the privacy and integrity of their data, with confidentiality enhanced by the use of encryption. Because speed and availability are limited by a user's online connection, users with large amounts of data may need to use cloud seeding and large-scale recovery.
Management
Various methods can be used to manage backup media, striking a balance between accessibility, security and cost. These media management methods are not mutually exclusive and are frequently combined to meet the user's needs. Using on-line disks for staging data before it is sent to a near-line tape library is a common example.
Online
Online backup storage is typically the most accessible type of data storage, and can begin a restore in milliseconds. An internal hard disk or a disk array (maybe connected to SAN) is an example of an online backup. This type of storage is convenient and speedy, but is vulnerable to being deleted or overwritten, either by accident, by malevolent action, or in the wake of a data-deleting virus payload.
Near-line
Nearline storage is typically less accessible and less expensive than online storage, but still useful for backup data storage. A mechanical device is usually used to move media units from storage into a drive where the data can be read or written. Generally it has safety properties similar to on-line storage. An example is a tape library with restore times ranging from seconds to a few minutes.
Off-line
Off-line storage requires some direct action to provide access to the storage media: for example, inserting a tape into a tape drive or plugging in a cable. Because the data is not accessible via any computer except during limited periods in which they are written or read back, they are largely immune to on-line backup failure modes. Access time varies depending on whether the media are on-site or off-site.
Off-site data protection
Backup media may be sent to an off-site vault to protect against a disaster or other site-specific problem. The vault can be as simple as a system administrator's home office or as sophisticated as a disaster-hardened, temperature-controlled, high-security bunker with facilities for backup media storage. A data replica can be off-site but also on-line (e.g., an off-site RAID mirror).
Backup site
A backup site or disaster recovery center is used to store data that can enable computer systems and networks to be restored and properly configured in the event of a disaster. Some organisations have their own data recovery centres, while others contract this out to a third-party. Due to high costs, backing up is rarely considered the preferred method of moving data to a DR site. A more typical way would be remote disk mirroring, which keeps the DR data as up to date as possible.
Selection and extraction of data
A backup operation starts with selecting and extracting coherent units of data. Most data on modern computer systems is stored in discrete units, known as files. These files are organized into filesystems. Deciding what to back up at any given time involves tradeoffs. By backing up too much redundant data, the information repository will fill up too quickly. Backing up an insufficient amount of data can eventually lead to the loss of critical information.
Files
Copying files: Making copies of files is the simplest and most common way to perform a backup. A means to perform this basic function is included in all backup software and all operating systems.
Partial file copying: A backup may include only the blocks or bytes within a file that have changed in a given period of time. This can substantially reduce needed storage space, but requires higher sophistication to reconstruct files in a restore situation. Some implementations require integration with the source file system.
Deleted files: To prevent the unintentional restoration of files that have been intentionally deleted, a record of the deletion must be kept.
Versioning of files: Most backup applications, other than those that do only full only/System imaging, also back up files that have been modified since the last backup. "That way, you can retrieve many different versions of a given file, and if you delete it on your hard disk, you can still find it in your [information repository] archive."
Filesystems
Filesystem dump: A copy of the whole filesystem in block-level can be made. This is also known as a "raw partition backup" and is related to disk imaging. The process usually involves unmounting the filesystem and running a program like dd (Unix). Because the disk is read sequentially and with large buffers, this type of backup can be faster than reading every file normally, especially when the filesystem contains many small files, is highly fragmented, or is nearly full. But because this method also reads the free disk blocks that contain no useful data, this method can also be slower than conventional reading, especially when the filesystem is nearly empty. Some filesystems, such as XFS, provide a "dump" utility that reads the disk sequentially for high performance while skipping unused sections. The corresponding restore utility can selectively restore individual files or the entire volume at the operator's choice.
Identification of changes: Some filesystems have an archive bit for each file that says it was recently changed. Some backup software looks at the date of the file and compares it with the last backup to determine whether the file was changed.
Versioning file system: A versioning filesystem tracks all changes to a file. The NILFS versioning filesystem for Linux is an example.
Live data
Files that are actively being updated present a challenge to back up. One way to back up live data is to temporarily quiesce them (e.g., close all files), take a "snapshot", and then resume live operations. At this point the snapshot can be backed up through normal methods. A snapshot is an instantaneous function of some filesystems that presents a copy of the filesystem as if it were frozen at a specific point in time, often by a copy-on-write mechanism. Snapshotting a file while it is being changed results in a corrupted file that is unusable. This is also the case across interrelated files, as may be found in a conventional database or in applications such as Microsoft Exchange Server. The term fuzzy backup can be used to describe a backup of live data that looks like it ran correctly, but does not represent the state of the data at a single point in time.
Backup options for data files that cannot be or are not quiesced include:
Open file backup: Many backup software applications undertake to back up open files in an internally consistent state. Some applications simply check whether open files are in use and try again later. Other applications exclude open files that are updated very frequently. Some low-availability interactive applications can be backed up via natural/induced pausing.
Interrelated database files backup: Some interrelated database file systems offer a means to generate a "hot backup" of the database while it is online and usable. This may include a snapshot of the data files plus a snapshotted log of changes made while the backup is running. Upon a restore, the changes in the log files are applied to bring the copy of the database up to the point in time at which the initial backup ended. Other low-availability interactive applications can be backed up via coordinated snapshots. However, genuinely-high-availability interactive applications can be only be backed up via Continuous Data Protection.
Metadata
Not all information stored on the computer is stored in files. Accurately recovering a complete system from scratch requires keeping track of this non-file data too.
System description: System specifications are needed to procure an exact replacement after a disaster.
Boot sector: The boot sector can sometimes be recreated more easily than saving it. It usually isn't a normal file and the system won't boot without it.
Partition layout: The layout of the original disk, as well as partition tables and filesystem settings, is needed to properly recreate the original system.
File metadata: Each file's permissions, owner, group, ACLs, and any other metadata need to be backed up for a restore to properly recreate the original environment.
System metadata: Different operating systems have different ways of storing configuration information. Microsoft Windows keeps a registry of system information that is more difficult to restore than a typical file.
Manipulation of data and dataset optimization
It is frequently useful or required to manipulate the data being backed up to optimize the backup process. These manipulations can improve backup speed, restore speed, data security, media usage and/or reduced bandwidth requirements.
Automated data grooming
Out-of-date data can be automatically deleted, but for personal backup applications—as opposed to enterprise client-server backup applications where automated data "grooming" can be customized—the deletion can at most be globally delayed or be disabled.
Compression
Various schemes can be employed to shrink the size of the source data to be stored so that it uses less storage space. Compression is frequently a built-in feature of tape drive hardware.
Deduplication
Redundancy due to backing up similarly configured workstations can be reduced, thus storing just one copy. This technique can be applied at the file or raw block level. This potentially large reduction is called deduplication. It can occur on a server before any data moves to backup media, sometimes referred to as source/client side deduplication. This approach also reduces bandwidth required to send backup data to its target media. The process can also occur at the target storage device, sometimes referred to as inline or back-end deduplication.
Duplication
Sometimes backups are duplicated to a second set of storage media. This can be done to rearrange the archive files to optimize restore speed, or to have a second copy at a different location or on a different storage medium—as in the disk-to-disk-to-tape capability of Enterprise client-server backup.
Encryption
High-capacity removable storage media such as backup tapes present a data security risk if they are lost or stolen. Encrypting the data on these media can mitigate this problem, however encryption is a CPU intensive process that can slow down backup speeds, and the security of the encrypted backups is only as effective as the security of the key management policy.
Multiplexing
When there are many more computers to be backed up than there are destination storage devices, the ability to use a single storage device with several simultaneous backups can be useful. However cramming the scheduled backup window via "multiplexed backup" is only used for tape destinations.
Refactoring
The process of rearranging the sets of backups in an archive file is known as refactoring. For example, if a backup system uses a single tape each day to store the incremental backups for all the protected computers, restoring one of the computers could require many tapes. Refactoring could be used to consolidate all the backups for a single computer onto a single tape, creating a "synthetic full backup". This is especially useful for backup systems that do incrementals forever style backups.
Staging
Sometimes backups are copied to a staging disk before being copied to tape. This process is sometimes referred to as D2D2T, an acronym for Disk-to-disk-to-tape. It can be useful if there is a problem matching the speed of the final destination device with the source device, as is frequently faced in network-based backup systems. It can also serve as a centralized location for applying other data manipulation techniques.
Objectives
Recovery point objective (RPO): The point in time that the restarted infrastructure will reflect, expressed as "the maximum targeted period in which data (transactions) might be lost from an IT service due to a major incident". Essentially, this is the roll-back that will be experienced as a result of the recovery. The most desirable RPO would be the point just prior to the data loss event. Making a more recent recovery point achievable requires increasing the frequency of synchronization between the source data and the backup repository.
Recovery time objective (RTO): The amount of time elapsed between disaster and restoration of business functions.
Data security: In addition to preserving access to data for its owners, data must be restricted from unauthorized access. Backups must be performed in a manner that does not compromise the original owner's undertaking. This can be achieved with data encryption and proper media handling policies.
Data retention period: Regulations and policy can lead to situations where backups are expected to be retained for a particular period, but not any further. Retaining backups after this period can lead to unwanted liability and sub-optimal use of storage media.
Checksum or hash function validation: Applications that back up to tape archive files need this option to verify that the data was accurately copied.
Backup process monitoring: Enterprise client-server backup applications need a user interface that allows administrators to monitor the backup process, and proves compliance to regulatory bodies outside the organization; for example, an insurance company in the USA might be required under HIPAA to demonstrate that its client data meet records retention requirements.
User-initiated backups and restores: To avoid or recover from minor disasters, such as inadvertently deleting or overwriting the "good" versions of one or more files, the computer user—rather than an administrator—may initiate backups and restores (from not necessarily the most-recent backup) of files or folders.
| Technology | Data storage and memory | null |
534057 | https://en.wikipedia.org/wiki/Triple%20bond | Triple bond | A triple bond in chemistry is a chemical bond between two atoms involving six bonding electrons instead of the usual two in a covalent single bond. Triple bonds are stronger than the equivalent single bonds or double bonds, with a bond order of three. The most common triple bond is in a nitrogen N2 molecule; the second most common is that between two carbon atoms, which can be found in alkynes. Other functional groups containing a triple bond are cyanides and isocyanides. Some diatomic molecules, such as diphosphorus and carbon monoxide, are also triple bonded. In skeletal formulae the triple bond is drawn as three parallel lines (≡) between the two connected atoms.
Bonding
Triple bonding can be explained in terms of orbital hybridization. In the case of acetylene, each carbon atom has two sp-orbitals and two p-orbitals. The two sp-orbitals are linear, with 180° bond angles, and occupy the x-axis in the cartesian coordinate system. The p-orbitals are perpendicular to the sp-orbitals on the y-axis and the z-axis. When the atoms approach each other, the sp orbitals overlap to form an sp-sp sigma bond. At the same time the pz-orbitals approach and together they form a pz-pz pi-bond. Likewise, the other pair of py-orbitals form a py-py pi-bond. The result is formation of one sigma bond and two pi bonds.
In the bent bond model, the triple bond can also formed by the overlapping of three sp3 lobes without the need to invoke a pi-bond.
Triple bonds between elements heavier than oxygen
Many elements beyond oxygen can form triple bonds. These bonds are common in some transition metals. Hexa(tert-butoxy)ditungsten(III) and Hexa(tert-butoxy)dimolybdenum(III) are well known examples, in which the metal-metal bond distance is about 233 pm. Hexa(tert-butoxy)ditungsten(III) has attracted particular attention for its reactions with alkynes, leading to metal-carbon triple bonded compounds of the formula RC≡W(OBut)3
Additionally, phosphorus can exist as the highly reactive diatomic molecule diphosphorus, which has roughly half the bond-dissociation energy of dinitrogen.
| Physical sciences | Bonding | Chemistry |
534313 | https://en.wikipedia.org/wiki/Long-tailed%20weasel | Long-tailed weasel | The long-tailed weasel (Neogale frenata), also known as the bridled weasel, masked ermine, or big stoat, is a species of weasel found in North, Central, and South America. It is distinct from the short-tailed weasel (Mustela erminea), also known as a "stoat", a close relation in the genus Mustela that originated in Eurasia and crossed into North America some half million years ago; the two species are visually similar, especially the black tail tip.
Long-tailed weasels exhibit scale-dependent patterns of habitat selection, favoring forest patches, fencerows, and drainage ditches while avoiding agricultural fields, suggesting sensitivity to habitat fragmentation due to agricultural practices.
Taxonomy
The long-tailed weasel was originally described in the genus Mustela with the name Mustela frenata by Hinrich Lichtenstein in 1831. In 1993, the classification, Mustela frenata, was accepted into the second edition of the Mammal species of the world: a taxonomic and geographic reference, which was published by the Smithsonian Institution Press. The species, with classification and name Mustela frenata, was accepted into the Global Biodiversity Information Facility. Later, in a study published in 2021 in the Journal of Animal Diversity, Bruce Patterson et al. reclassified the long-tailed weasel into the genus Neogale along with 2 other former Mustela species, as well as the two species formerly classified in Neovison.
Evolution
The long-tailed weasel is the product of a process begun 5–7 million years ago, when northern forests were replaced by open grassland, thus prompting an explosive evolution of small, burrowing rodents. The long-tailed weasel's ancestors were larger than the current form, and underwent a reduction in size to exploit the new food source. The long-tailed weasel arose in North America 2 million years ago, shortly before the stoat evolved as its mirror image in Eurasia. The species thrived during the Ice Age, as its small size and long body allowed it to easily operate beneath snow, as well as hunt in burrows. The long-tailed weasel and the stoat remained separated until half a million years ago, when falling sea levels exposed the Bering land bridge, thus allowing the stoat to cross into North America. However, unlike the latter species, the long-tailed weasel never crossed the land bridge, and did not spread into Eurasia.
Description
The long-tailed weasel is one of the larger weasels (comprising both Neogale and Mustela) in North America. There is substantial disagreement both on the upper end of their size and difference in size by sex by source: one indicates a body length of and a tail comprising 40–70% of the head and body length. It adds that in most populations, females are 10–15% smaller than males, thus making them about the same size as large male stoats, according to a second source. A third states they range from 11 to 22 inches (280–560 mm) in length, with the tail measuring an additional 3 to 6 inches (80–150 mm). It maintains the long-tailed weasel weighs between 3 and 9 ounces (85-267 g) with males being about twice as large as the females.
The eyes are black in daylight, but glow bright emerald green when caught in a spotlight at night. The dorsal fur is brown in summer, while the underparts are whitish and tinged with yellowish or buffy brown from the chin to the inguinal region. The tail has a distinct black tip. Long-tailed weasels in Florida and the southwestern US may have facial markings of a white or yellowish colour. In northern areas in winter, the long-tailed weasel's fur becomes white, sometimes with yellow tints, but the tail retains its black tip. The long-tailed weasel moults twice annually, once in autumn (October to mid-November) and once in spring (March–April). Each moult takes about 3–4 weeks and is governed by day length and mediated by the pituitary gland. Unlike the stoat, whose soles are thickly furred all year, the long-tailed weasel's soles are naked in summer. The long-tailed weasel has well-developed anal scent glands, which produce a strong and musky odour. Analysis of a dichloromethane extract of the anal gland secretion showed it contained 2,2-dimethylthietane, 2,4-dimethylthietane, 2,3-dimethylthietane, 2-propylthietane, 3,3-dimethyl-1,2-dithiolane, 3-ethyl-1,2-dithiolane, indole and 2-aminoacetophenone. Unlike skunks, which spray their musk, the long-tailed weasel drags and rubs its body over surfaces in order to leave the scent, to mark their territory and, when startled or threatened, to discourage predators.
Identification
Tracks and scat
The footprint of a long-tailed weasel is about 1 inch (25 mm) long. Although they have five toes, only four of them can be seen in their tracks. The only exception to this is when walking in the snow or mud, all five of their toes are shown. Their footprints will also appear heavier if the weasel is carrying food. Another way to determine the presence of a weasel is by looking for wavy indents made by their tails in the snow.
The long-tailed weasel uses one spot to leave their feces. This spot is usually near where they burrow. They'll continuously use this spot for their droppings until it gets covered by environmental changes.
Distinguishing features
A black-tipped tail, yellowish-white belly fur, and brown fur on its back and sides are distinguishing for the long-tailed weasel. Additionally, the long-tailed weasel has long whiskers, a long narrow body, and a long tail that is approximately half the length of the body and head of the weasel. Compared to the short-tailed weasel the long-tailed weasel lacks a white line on the insides of its legs.
Behaviour
Reproduction and development
The long-tailed weasel mates in July–August, with implantation of the fertilized egg on the uterine wall being delayed until about March. The gestation period lasts 10 months, with actual embryonic development taking place only during the last four weeks of this period, an adaptation to timing births for spring, when small mammals are abundant. Litter size generally consists of 5–8 kits, which are born in April–May. The kits are born partially naked, blind and weighing , about the same weight of a hummingbird. The long-tailed weasel's growth rate is rapid, as by the age of three weeks, the kits are well furred, can crawl outside the nest and eat meat. At this time, the kits weigh . At five weeks of age, the kit's eyes open, and the young become physically active and vocal. Weaning begins at this stage, with the kits emerging from the nest and accompanying the mother in hunting trips a week later. The kits are fully grown by autumn, at which time the family disbands. The females are able to breed at 3–4 months of age, while males become sexually mature at 15–18 months.
Denning and sheltering behaviour
The long-tailed weasel dens in ground burrows, under stumps or beneath rock piles. It usually does not dig its own burrows, but commonly uses abandoned chipmunk holes. The diameter nest chamber is situated around from the burrow entrance, and is lined with straw and the fur of prey.
Defense
The enemies of the long-tailed weasel are usually coyotes, foxes, wildcats, wolves, and the Canadian lynx. The weasel will give off its musky odor, however, this is not primarily used when encountering other creatures. When leaving an area they were just in, they will leave their odor behind. This is done by the weasels taking themselves and hauling their bodies across surfaces they just interacted with. The long-tailed weasel does this to "discourage predators" from coming back to the area, possibly indicating that the weasel considers this a safe haven for return. This type of reaction is also reserved for when the weasel feels it is in danger, or when it is looking for a mate. Tree-climbing is another type of defense mechanism that long-tailed weasels utilize against predators on the ground. These weasels will climb up a reasonable height of a tree when they sense that they are in danger. They will then sit silent and "motionless", while looking at their presumed predator. These weasels keep their guard up like this until the predator leaves, and when the weasel considers itself no longer in danger.
Another common defense of long-tailed weasels is its black-tipped tail, which differs in color from the rest of the body. When the long-tailed weasel becomes more white in the winter, this defense mechanism is especially used. The black-tipped tail distracts predators from the rest of the body, as it is more visible to the eye of a predator. This causes the visibility of the actual weasel to be rather difficult and makes the predator attack the tail instead of the weasel. The weasel is allowed to escape the predator because of this.
Diet
The long-tailed weasel is a fearless and aggressive hunter which may attack animals far larger than itself. When stalking, it waves its head from side to side in order to pick up the scent of its prey. It hunts small prey, such as mice, by rushing at them and killing them with one bite to the head. With large prey, such as rabbits, the long-tailed weasel strikes quickly, taking its prey off guard. It grabs the nearest part of the animal and climbs upon its body, maintaining its hold with its feet. The long-tailed weasel then manoeuvres itself to inflict a lethal bite to the neck.
The long-tailed weasel is an obligate carnivore which prefers its prey to be fresh or alive, eating only the carrion stored within its burrows. Rodents are almost exclusively taken when they are available. Its primary prey consists of mice, rats, squirrels, chipmunks, shrews, moles and rabbits. Occasionally, it may eat small birds, bird eggs, reptiles, amphibians, fish, earthworms and some insects. The species has also been observed to take bats from nursery colonies. It occasionally surplus kills, usually in spring when the kits are being fed, and again in autumn. Some of the surplus kills may be cached, but are usually left uneaten. Kits in captivity eat from a quarter to half of their body weight in 24 hours, while adults eat only one fifth to one third. After killing its prey, the long-tailed weasel laps up the blood, but does not suck it, as is popularly believed. With small prey, also the fur, feathers, flesh and bones are consumed, but only some flesh is eaten from large prey. When stealing eggs, the long-tailed weasel removes each egg from its nest one at a time, then carries it in its mouth to a safe location where it bites off the top and licks out the contents or if they have babies in the den they may hold it in their mouth all the way back to them.
Subspecies
, 42 subspecies are recognised.
Cultural meanings
In North America, Native Americans (in the region of Chatham County, North Carolina) deemed the long-tailed weasel to be a bad sign; crossing its path meant a "speedy death".
| Biology and health sciences | Mustelidae | Animals |
534435 | https://en.wikipedia.org/wiki/Lateral%20line | Lateral line | The lateral line, also called the lateral line organ (LLO), is a system of sensory organs found in fish, used to detect movement, vibration, and pressure gradients in the surrounding water. The sensory ability is achieved via modified epithelial cells, known as hair cells, which respond to displacement caused by motion and transduce these signals into electrical impulses via excitatory synapses. Lateral lines play an important role in schooling behavior, predation, and orientation.
Early in the evolution of fish, some of the sensory organs of the lateral line were modified to function as the electroreceptors called ampullae of Lorenzini. The lateral line system is ancient and basal to the vertebrate clade, as it is found in fishes that diverged over 400 million years ago.
Function
The lateral line system allows the detection of movement, vibration, and pressure gradients in the water surrounding an animal. It plays an essential role in orientation, predation, and fish schooling by providing spatial awareness and the ability to navigate in the environment. Analysis has shown that the lateral line system should be an effective passive sensing system able to discriminate between submerged obstacles by their shape. The lateral line allows fish to navigate and hunt in water with poor visibility.
The lateral line system enables predatory fishes to detect vibrations made by their prey, and to orient towards the source to begin predatory action. Blinded predatory fishes remain able to hunt, but not when lateral line function is inhibited by cobalt ions.
The lateral line plays a role in fish schooling. Blinded Pollachius virens were able to integrate into a school, whereas fish with severed lateral lines could not. It may have evolved further to allow fish to forage in dark caves. In Mexican blind cave fish, Astyanax mexicanus, neuromasts in and around the orbit of the eye are bigger and around twice as sensitive as those of surface-living fish.
One function of schooling may be to confuse the lateral line of predatory fishes. A single prey fish creates a simple particle velocity pattern, whereas the pressure gradients of many closely swimming (schooling) prey fish overlap, creating a complex pattern. This makes it difficult for predatory fishes to identify individual prey through lateral line perception.
Anatomy
Lateral lines are usually visible as faint lines of pores running along each side of a fish's body. The functional units of the lateral line are the neuromasts, discrete mechanoreceptive organs that sense movement in water. There are two main varieties: canal neuromasts and superficial neuromasts. Superficial neuromasts are on the surface of the body, while canal neuromasts are along the lateral lines in subdermal, fluid-filled canals. Each neuromast consists of receptive hair cells whose tips are covered by a flexible jellylike cupula. Hair cells typically possess both glutamatergic afferent connections and cholinergic efferent connections. The receptive hair cells are modified epithelial cells; they typically possess bundles of 40–50 microvilli "hairs" which function as the mechanoreceptors. Within each bundle, the hairs are organized in a rough "staircase" from shortest to longest.
Signal transduction
The hair cells are stimulated by the deflection of their hair bundles in the direction of the tallest "hairs" or stereocilia. The deflection allows cations to enter through a mechanically gated channel, causing depolarization or hyperpolarization of the hair cell. Depolarization opens Cav1.3 calcium channels in the basolateral membrane.
Hair cells use a system of transduction with rate coding to transmit the directionality of a stimulus. The hair cells produce a constant, tonic rate of firing. As mechanical motion is transmitted through water to the neuromast, the cupula bends and is displaced according to the strength of the stimulus. This results in a shift in the cell's ionic permeability. Deflection towards the longest hair results in depolarization of the hair cell, increased neurotransmitter release at the excitatory afferent synapse, and a higher rate of signal transduction. Deflection towards the shorter hair has the opposite effect, hyperpolarizing the hair cell and producing a decreased rate of neurotransmitter release. These electrical impulses are then transmitted along afferent lateral neurons to the brain.
While both varieties of neuromasts utilize this method of transduction, their specialized organization gives them different mechanoreceptive capacities. Superficial organs are exposed more directly to the external environment. The organization of the bundles within their organs is seemingly haphazard, incorporating various shapes and sizes of microvilli within bundles. This suggests coarse but wide-ranging detection. In contrast, the structure of canal organs allow canal neuromasts more sophisticated mechanoreception, such as of pressure differentials. As current moves across the pores, a pressure differential is created, inducing a flow in the canal fluid. This moves the cupulae of the neuromasts in the canal, resulting in a deflection of the hairs in the direction of the flow.
Electrophysiology
The mechanoreceptive hair cells of the lateral line structure are integrated into more complex circuits through their afferent and efferent connections. The synapses that directly participate in the transduction of mechanical information are excitatory afferent connections that utilize glutamate. Species vary in their neuromast and afferent connections, providing differing mechanoreceptive properties. For instance, the superficial neuromasts of the midshipman fish, Porichthys notatus, are sensitive to specific stimulation frequencies. One variety is attuned to collect information about acceleration, at stimulation frequencies between 30 and 200 Hz. The other type obtains information about velocity, and is most receptive to stimulation below 30 Hz.
The efferent synapses to hair cells are inhibitory and use acetylcholine as a transmitter. They are crucial participants in a corollary discharge system designed to limit self-generated interference. When a fish moves, it creates disturbances in the water that could be detected by the lateral line system, potentially interfering with the detection of other biologically relevant signals. To prevent this, an efferent signal is sent to the hair cell upon motor action, resulting in inhibition which counteracts the excitation resulting from reception of the self-generated stimulation. This allows the fish to detect external stimuli without interference from its own movements.
Signals from the hair cells are transmitted along lateral neurons to the brain. The area where these signals most often terminate is the medial octavolateralis nucleus (MON), which probably processes and integrates mechanoreceptive information. The deep MON contains distinct layers of basilar and non-basilar crest cells, suggesting computational pathways analogous to the electrosensory lateral line lobe of electric fish. The MON is likely involved in the integration of excitatory and inhibitory parallel circuits to interpret mechanoreceptive information.
Evolution
The use of mechanosensitive hairs is homologous to the functioning of hair cells in the auditory and vestibular systems, indicating a close link between these systems. Due to many overlapping functions and their great similarity in ultrastructure and development, the lateral line system and the inner ear of fish are often grouped together as the octavolateralis system (OLS). Here, the lateral line system detects particle velocities and accelerations with frequencies below 100 Hz. These low frequencies create large wavelengths, which induce strong particle accelerations in the near field of swimming fish that do not radiate into the far field as acoustic waves due to an acoustic short circuit. The auditory system detects pressure fluctuations with frequencies above 100 Hz that propagate to the far field as waves.
The lateral line system is ancient and basal to the vertebrate clade; it is found in groups of fishes that diverged over 400 million years ago, including the lampreys, cartilaginous fishes, and bony fishes. Most amphibian larvae and some fully aquatic adult amphibians possess mechanosensitive systems comparable to the lateral line. The terrestrial tetrapods have secondarily lost their lateral line organs, which are ineffective when not submerged.
The electroreceptive organs, called ampullae of Lorenzini, appearing as pits in the skin of sharks and some other fishes, evolved from the lateral line organ. Passive electroreception using ampullae is an ancestral trait in the vertebrates, meaning that it was present in their last common ancestor.
| Biology and health sciences | Sensory nervous system | Biology |
534710 | https://en.wikipedia.org/wiki/Cone%20cell | Cone cell | Cone cells or cones are photoreceptor cells in the retinas of vertebrates' eyes. They respond differently to different wavelengths of light, and the combination of their responses is responsible for color vision. Cones function best in relatively bright light, called the photopic region, as opposed to rod cells, which work better in dim light, or the scotopic region. Cone cells are densely packed in the fovea centralis, a 0.3 mm diameter rod-free area with very thin, densely packed cones which quickly reduce in number towards the periphery of the retina. Conversely, like rods, they are absent from the optic disc, contributing to the blind spot. There are about six to seven million cones in a human eye (vs ~92 million rods), with the highest concentration being towards the macula.
Cones are less sensitive to light than the rod cells in the retina (which support vision at low light levels), but allow the perception of color. They are also able to perceive finer detail and more rapid changes in images because their response times to stimuli are faster than those of rods. In humans, cones are normally one of three types: S-cones, M-cones and L-cones, with each type bearing a different opsin: OPN1SW, OPN1MW, and OPN1LW respectively. These cones are sensitive to visible wavelengths of light that correspond to short-wavelength, medium-wavelength and longer-wavelength light respectively. Because humans usually have three kinds of cones with different photopsins, which have different response curves and thus respond to variation in color in different ways, humans have trichromatic vision. Being color blind can change this, and there have been some verified reports of people with four types of cones, giving them tetrachromatic vision.
The three pigments responsible for detecting light have been shown to vary in their exact chemical composition due to genetic mutation; different individuals will have cones with different color sensitivity.
Structure
Types
Humans normally have three types of cones, usually designated L, M and S for long, medium and short wavelengths respectively. L cones respond most strongly to light of the longer red wavelengths, peaking at about 560 nm, and make up the majority of the cones in the human eye. M cones, the second most common type, respond most strongly to yellow to green medium-wavelength light, peaking at 530 nm, and make up about a third of cones in the human eye. S cones respond most strongly to blue short-wavelength light, peaking at 420 nm, and make up only around 2% of the cones in the human retina. The peak wavelengths of L, M, and S cones occur in the ranges of 564–580 nm, 534–545 nm, and 420–440 nm, respectively, depending on the individual. Differences between the three types of cone cells are due to the distinct opsins they express, which are coded for by the genes OPN1LW, OPN1MW, and OPN1SW, respectively. The CIE 1931 color space is an often-used model of spectral sensitivities of the three cells of an average human.
Some non-human animals have different numbers of cone types (see Color_vision#Color_vision_in_nonhumans).
Shape and arrangement
Cone cells are shorter but wider than rod cells. While rods outnumber cones in most parts of the retina, the fovea, responsible for sharp central vision, consists almost entirely of cones. Structurally, cone cells have a cone-like shape at one end where an opsin specific to the type of cone absorbs incoming light. They are typically 40–50μm long, and their diameter varies from 0.5 to 4.0μm, being smallest and most tightly packed at the center of the eye at the fovea. The S cone spacing is slightly larger than the others.
Photobleaching can be used to determine cone arrangement. This is done by exposing dark-adapted retina to a certain wavelength of light that paralyzes the particular type of cone sensitive to that wavelength for up to thirty minutes from being able to dark-adapt, making it appear white in contrast to the grey dark-adapted cones when a picture of the retina is taken. The results illustrate that S cones are randomly placed and appear much less frequently than the M and L cones. The ratio of M and L cones varies greatly among different people with regular vision (e.g. values of 75.8% L with 20.0% M versus 50.6% L with 44.2% M in two male subjects).
Like rods, each cone cell has a synaptic terminal, inner and outer segments, as well as an interior nucleus and various mitochondria. The synaptic terminal forms a synapse with a neuron bipolar cell. The inner and outer segments are connected by a cilium. The inner segment contains organelles and the cell's nucleus, while the outer segment contains the light-absorbing materials.
The outer segments of cones have invaginations of their cell membranes that create stacks of membranous disks. Photopigments exist as transmembrane proteins within these disks, which provide more surface area for light to affect the pigments. In cones, these disks are attached to the outer membrane, whereas they are pinched off and exist separately in rods. Neither rods nor cones divide, but their membranous disks wear out and are worn off at the end of the outer segment, to be consumed and recycled by phagocytic cells.
Function
The difference in the signals received from the three cone types allows the brain to perceive a continuous range of colors through the opponent process of color vision. Rod cells have a peak sensitivity at 498 nm, roughly halfway between the peak sensitivities of the S and M cones.
All of the receptors contain the protein photopsin. Variations in its conformation cause differences in the optimum wavelengths absorbed.
The color yellow, for example, is perceived when the L cones are stimulated slightly more than the M cones, and the color red is perceived when the L cones are stimulated significantly more than the M cones. Similarly, blue and violet hues are perceived when the S receptor is stimulated more. S Cones are most sensitive to light at wavelengths around 420 nm. However, the lens and cornea of the human eye are increasingly absorptive to shorter wavelengths, and this sets the short wavelength limit of human-visible light to approximately 380 nm, which is therefore called 'ultraviolet' light. People with aphakia, a condition where the eye lacks a lens, sometimes report the ability to see into the ultraviolet range. At moderate to bright light levels where the cones function, the eye is more sensitive to yellowish-green light than other colors because this stimulates the two most common (M and L) of the three kinds of cones almost equally. At lower light levels, where only the rod cells function, the sensitivity is greatest at a blueish-green wavelength.
Cones also tend to possess a significantly elevated visual acuity because each cone cell has a lone connection to the optic nerve, therefore, the cones have an easier time telling that two stimuli are isolated. Separate connectivity is established in the
inner plexiform layer so that each connection is parallel.
The response of cone cells to light is also directionally nonuniform, peaking at a direction that receives light from the center of the pupil; this effect is known as the Stiles–Crawford effect.
S cones may play a role in the regulation of the circadian system and the secretion of melatonin but this role is not clear yet. The exact contribution of S cone activation to circadian regulation is unclear but any potential role would be secondary to the better established role of melanopsin (see also Intrinsically photosensitive retinal ganglion cell).
Color afterimage
Sensitivity to a prolonged stimulation tends to decline over time, leading to neural adaptation. An interesting effect occurs when staring at a particular color for a minute or so. Such action leads to an exhaustion of the cone cells that respond to that color – resulting in the afterimage. This vivid color aftereffect can last for a minute or more.
Associated diseases
Achromatopsia (rod monochromacy)a form of monochromacy with no functional cones
Blue cone monochromacya rare form of monochromacy with only functional S-cones
Congenital red–green color blindnesspartial color blindness include protanopia, deuteranopia, etc.
Oligocone trichromacypoor visual acuity and impairment of cone function according to ERG, but without significant color vision loss.
Bradyopsiaphotopic vision cannot respond quickly to stimuli.
Bornholm eye diseaseX-linked recessive myopia, astigmatism, impaired visual acuity and red–green dichromacy.
Cone dystrophya degenerative loss of cone cells
Retinoblastomaa type of cancer originating from cone precursor cells
| Biology and health sciences | Visual system | Biology |
534914 | https://en.wikipedia.org/wiki/Spectral%20graph%20theory | Spectral graph theory | In mathematics, spectral graph theory is the study of the properties of a graph in relationship to the characteristic polynomial, eigenvalues, and eigenvectors of matrices associated with the graph, such as its adjacency matrix or Laplacian matrix.
The adjacency matrix of a simple undirected graph is a real symmetric matrix and is therefore orthogonally diagonalizable; its eigenvalues are real algebraic integers.
While the adjacency matrix depends on the vertex labeling, its spectrum is a graph invariant, although not a complete one.
Spectral graph theory is also concerned with graph parameters that are defined via multiplicities of eigenvalues of matrices associated to the graph, such as the Colin de Verdière number.
Cospectral graphs
Two graphs are called cospectral or isospectral if the adjacency matrices of the graphs are isospectral, that is, if the adjacency matrices have equal multisets of eigenvalues.
Cospectral graphs need not be isomorphic, but isomorphic graphs are always cospectral.
Graphs determined by their spectrum
A graph is said to be determined by its spectrum if any other graph with the same spectrum as is isomorphic to .
Some first examples of families of graphs that are determined by their spectrum include:
The complete graphs.
The finite starlike trees.
Cospectral mates
A pair of graphs are said to be cospectral mates if they have the same spectrum, but are non-isomorphic.
The smallest pair of cospectral mates is {K1,4, C4 ∪ K1}, comprising the 5-vertex star and the graph union of the 4-vertex cycle and the single-vertex graph. The first example of cospectral graphs was reported by Collatz and Sinogowitz in 1957.
The smallest pair of polyhedral cospectral mates are enneahedra with eight vertices each.
Finding cospectral graphs
Almost all trees are cospectral, i.e., as the number of vertices grows, the fraction of trees for which there exists a cospectral tree goes to 1.
A pair of regular graphs are cospectral if and only if their complements are cospectral.
A pair of distance-regular graphs are cospectral if and only if they have the same intersection array.
Cospectral graphs can also be constructed by means of the Sunada method.
Another important source of cospectral graphs are the point-collinearity graphs and the line-intersection graphs of point-line geometries. These graphs are always cospectral but are often non-isomorphic.
Cheeger inequality
The famous Cheeger's inequality from Riemannian geometry has a discrete analogue involving the Laplacian matrix; this is perhaps the most important theorem in spectral graph theory and one of the most useful facts in algorithmic applications. It approximates the sparsest cut of a graph through the second eigenvalue of its Laplacian.
Cheeger constant
The Cheeger constant (also Cheeger number or isoperimetric number) of a graph is a numerical measure of whether or not a graph has a "bottleneck". The Cheeger constant as a measure of "bottleneckedness" is of great interest in many areas: for example, constructing well-connected networks of computers, card shuffling, and low-dimensional topology (in particular, the study of hyperbolic 3-manifolds).
More formally, the Cheeger constant h(G) of a graph G on n vertices is defined as
where the minimum is over all nonempty sets S of at most n/2 vertices and ∂(S) is the edge boundary of S, i.e., the set of edges with exactly one endpoint in S.
Cheeger inequality
When the graph G is d-regular, there is a relationship between h(G) and the spectral gap d − λ2 of G. An inequality due to Dodziuk and independently Alon and Milman states that
This inequality is closely related to the Cheeger bound for Markov chains and can be seen as a discrete version of Cheeger's inequality in Riemannian geometry.
For general connected graphs that are not necessarily regular, an alternative inequality is given by Chung
where is the least nontrivial eigenvalue of the normalized Laplacian, and is the (normalized) Cheeger constant
where is the sum of degrees of vertices in .
Hoffman–Delsarte inequality
There is an eigenvalue bound for independent sets in regular graphs, originally due to Alan J. Hoffman and Philippe Delsarte.
Suppose that is a -regular graph on vertices with least eigenvalue . Then:where denotes its independence number.
This bound has been applied to establish e.g. algebraic proofs of the Erdős–Ko–Rado theorem and its analogue for intersecting families of subspaces over finite fields.
For general graphs which are not necessarily regular, a similar upper bound for the independence number can be derived by using the maximum eigenvalue
of the normalized Laplacian of :
where and denote the maximum and minimum degree in , respectively. This a consequence of a more general inequality (pp. 109 in
):
where is an independent set of vertices and denotes the sum of degrees of vertices in .
Historical outline
Spectral graph theory emerged in the 1950s and 1960s. Besides graph theoretic research on the relationship between structural and spectral properties of graphs, another major source was research in quantum chemistry, but the connections between these two lines of work were not discovered until much later. The 1980 monograph Spectra of Graphs by Cvetković, Doob, and Sachs summarised nearly all research to date in the area. In 1988 it was updated by the survey Recent Results in the Theory of Graph Spectra. The 3rd edition of Spectra of Graphs (1995) contains a summary of the further recent contributions to the subject. Discrete geometric analysis created and developed by Toshikazu Sunada in the 2000s deals with spectral graph theory in terms of discrete Laplacians associated with weighted graphs, and finds application in various fields, including shape analysis. In most recent years, the spectral graph theory has expanded to vertex-varying graphs often encountered in many real-life applications.
| Mathematics | Graph theory | null |
535191 | https://en.wikipedia.org/wiki/Real-time%20clock | Real-time clock | A real-time clock (RTC) is an electronic device (most often in the form of an integrated circuit) that measures the passage of time.
Although the term often refers to the devices in personal computers, servers and embedded systems, RTCs are present in almost any electronic device which needs to keep accurate time of day.
Terminology
The term real-time clock is used to avoid confusion with ordinary hardware clocks which are only signals that govern digital electronics, and do not count time in human units. RTC should not be confused with real-time computing, which shares its three-letter acronym but does not directly relate to time of day.
Purpose
Although keeping time can be done without an RTC, using one has benefits:
Reliably maintains and provides current time through disruptive system states such as hangs, sleep, reboots, or if given sufficient backup power, full shutdown and hardware reassembly, without the need to have its time set again.
Low power consumption (important when running from alternate power)
Frees the main system for time-critical tasks
Sometimes more accurate than other methods
A GPS receiver can shorten its startup time by comparing the current time, according to its RTC, with the time at which it last had a valid signal. If it has been less than a few hours, then the previous ephemeris is still usable.
Some motherboards are made without RTCs. The RTC may be omitted out of desire to save money or reduce possible sources of hardware failure.
Power source
RTCs often have an alternate source of power, so they can continue to keep time while the primary source of power is off or unavailable. This alternate source of power is normally a lithium battery in older systems, but some newer systems use a supercapacitor, because they are rechargeable and can be soldered. The alternate power source can also supply power to battery backed RAM.
Timing
Most RTCs use a crystal oscillator, but some have the option of using the power line frequency. The crystal frequency is usually 32.768 kHz, the same frequency used in quartz clocks and watches. Being exactly 215 cycles per second, it is a convenient rate to use with simple binary counter circuits. The low frequency saves power, while remaining above human hearing range. The quartz tuning fork of these crystals does not change size much from temperature, so temperature does not change its frequency much.
Some RTCs use a micromechanical resonator on the silicon chip of the RTC. This reduces the size and cost of an RTC by reducing its parts count. Micromechanical resonators are much more sensitive to temperature than quartz resonators. So, these compensate for temperature changes using an electronic thermometer and electronic logic.
Typical crystal RTC accuracy specifications are from ±100 to ±20 parts per million (8.6 to 1.7 seconds per day), but temperature-compensated RTC ICs are available accurate to less than 5 parts per million. In practical terms, this is good enough to perform celestial navigation, the classic task of a chronometer. In 2011, chip-scale atomic clocks became available. Although vastly more expensive and power-hungry (120 mW vs. <1 μW), they keep time within 50 parts per trillion ().
Examples
Many integrated circuit manufacturers make RTCs, including Epson, Intersil, IDT, Maxim, NXP Semiconductors, Texas Instruments, STMicroelectronics and Ricoh. A common RTC used in single-board computers is the Maxim Integrated DS1307.
The RTC was introduced to PC compatibles by the IBM PC/AT in 1984, which used a Motorola MC146818 RTC. Later, Dallas Semiconductor made compatible RTCs, which were often used in older personal computers, and are easily found on motherboards because of their distinctive black battery cap and silkscreened logo. A standard CMOS interface is available for the PC RTC.
In newer computer systems, the RTC is integrated into the southbridge chip.
Some microcontrollers have a real-time clock built in, generally only the ones with many other features and peripherals.
Radio-based RTCs
Some modern computers receive clock information by digital radio and use it to promote time-standards. There are two common methods: Most cell phone protocols (e.g. LTE) directly provide the current local time. If an internet radio is available, a computer may use the network time protocol. Computers used as local time servers occasionally use GPS or ultra-low frequency radio transmissions broadcast by a national standards organization (i.e. a radio clock).
Software-based RTCs
The following system is well-known to embedded systems programmers, who sometimes must construct RTCs in systems that lack them. Most computers have one or more hardware timers that use timing signals from quartz crystals or ceramic resonators. These have inaccurate absolute timing (more than 100 parts per million) that is yet very repeatable (often less than 1 ppm). Software can do the math to make these into accurate RTCs. The hardware timer can produce a periodic interrupt, e.g. 50 Hz, to mimic a historic RTC (see below). However, it uses math to adjust the timing chain for accuracy:
time = time + rate.
When the "time" variable exceeds a constant, usually a power of two, the nominal, calculated clock time (say, for 1/50 of a second) is subtracted from "time", and the clock's timing-chain software is invoked to count fractions of seconds, seconds, etc. With 32-bit variables for time and rate, the mathematical resolution of "rate" can exceed one part per billion. The clock remains accurate because it will occasionally skip a fraction of a second, or increment by two fractions. The tiny skip ("jitter") is imperceptible for almost all real uses of an RTC.
The complexity with this system is determining the instantaneous corrected value for the variable "rate". The simplest system tracks RTC time and reference time between two settings of the clock, and divides reference time by RTC time to find "rate". Internet time is often accurate to less than 20 milliseconds, so 8000 or more seconds (2.2 or more hours) of separation between settings can usually divide the forty milliseconds (or less) of error to less than 5 parts per million to get chronometer-like accuracy. The main complexity with this system is converting dates and times to counts of seconds, but methods are well known.
If the RTC runs when a unit is off, usually the RTC will run at two rates, one when the unit is on and another when off. This is because the temperature and power-supply voltage in each state is consistent. To adjust for these states, the software calculates two rates. First, software records the RTC time, reference time, on seconds and off seconds for the two intervals between the last three times that the clock is set. Using this, it can measure the accuracy of the two intervals, with each interval having a different distribution of on and off seconds. The rate math solves two linear equations to calculate two rates, one for on and the other for off.
Another approach measures the temperature of the oscillator with an electronic thermometer, (e.g. a thermistor and analog-to-digital converter) and uses a polynomial to calculate "rate" about once per minute. These require a calibration that measures the frequency at several temperatures, and then a linear regression to find the equation of temperature. The most common quartz crystals in a system are SC-cut crystals, and their rates over temperature can be characterized with a 3rd-degree polynomial. So, to calibrate these, the frequency is measured at four temperatures. The common tuning-fork-style crystals used in watches and many RTC components have parabolic (2nd-degree) equations of temperature, and can be calibrated with only 3 measurements. MEMS oscillators vary, from 3rd degree to fifth degree polynomials, depending on their mechanical design, and so need from four to six calibration measurements. Something like this approach might be used in commercial RTC ICs, but the actual methods of efficient high-speed manufacturing are proprietary.
Historic RTCs
Some computer designs such as smaller IBM System/360s,
PDP-8s and Novas used a real-time clock that was accurate, simple and low cost. In Europe, North America and some other grids, the frequency of the AC mains is adjusted to the long-term frequency accuracy of the national standards. In those grids, clocks using AC mains can keep perfect time without adjustment. Such clocks are not practical in portable computers or grids (e.g. in South Asia) that do not regulate the frequency of AC mains.
These computers' power supplies use a transformer or resistor divider to produce a sine wave at logic voltages. This signal is conditioned by a zero crossing detector, either using a linear amplifier, or a schmitt trigger. The result is a square wave with single, fast edges at the mains frequency. This logic signal triggers an interrupt. The interrupt handler software usually counts cycles, seconds, etc. In this way, it can provide an entire clock and calendar. In the IBM 360, the interrupt updates a 64-bit count of microseconds utilized by standardized systems software. The clock's jitter error is half if the clock interrupts for each zero crossing, instead of each cycle.
The clock also usually formed the basis of computers' software timing chains; e.g. it was usually the timer used to switch tasks in an operating system. Counting timers used in modern computers provide similar features at lower precision, and may trace their requirements to this type of clock. (e.g. in the PDP-8, the mains-based clock, model DK8EA, came first, and was later followed by a crystal-based clock, DK8EC.)
A software-based clock must be set each time its computer is turned on. Originally this was done by computer operators. When the Internet became commonplace, network time protocols were used to automatically set clocks of this type.
| Technology | Computer hardware | null |
535458 | https://en.wikipedia.org/wiki/Stellar%20atmosphere | Stellar atmosphere | The stellar atmosphere is the outer region of the volume of a star, lying above the stellar core, radiation zone and convection zone.
Overview
The stellar atmosphere is divided into several regions of distinct character:
The photosphere, which is the atmosphere's lowest and coolest layer, is normally its only visible part. Light escaping from the surface of the star stems from this region and passes through the higher layers. The Sun's photosphere has a temperature in the range. Starspots, cool regions of disrupted magnetic field, lie in the photosphere.
Above the photosphere lies the chromosphere. This part of the atmosphere first cools down and then starts to heat up to about 10 times the temperature of the photosphere.
Above the chromosphere lies the transition region, where the temperature increases rapidly on a distance of only around .
Additionally, many stars have a molecular layer (MOLsphere) above the photosphere and just beyond or even within the chromosphere. The molecular layer is cool enough to contain molecules rather than plasma, and may consist of such components as carbon monoxide, water vapor, silicon monoxide, and titanium oxide.
The outermost part of the stellar atmosphere, or upper stellar atmosphere, is the corona, a tenuous plasma which has a temperature above one million Kelvin. While all stars on the main sequence feature transition regions and coronae, not all evolved stars do so. It seems that only some giants, and very few supergiants, possess coronae. An unresolved problem in stellar astrophysics is how the corona can be heated to such high temperatures. The answer is believed to lie in magnetic fields, but the exact mechanism remains unclear.
The astrosphere, which is in the case of the Sun the heliosphere, can be in a broader understanding considered the furthest part of a stellar atmosphere, before interstellar space begins at the heliopause. The astrosphere is not to be confused with the Solar System and its outermost region the Oort cloud, which extends much further than the astrosphere, therefore far into interstellar space.
During a total solar eclipse, the photosphere of the Sun is obscured, revealing its atmosphere's other layers. Observed during eclipse, the Sun's chromosphere appears (briefly) as a thin pinkish arc, and its corona is seen as a tufted halo. The same phenomenon in eclipsing binaries can make the chromosphere of giant stars visible.
| Physical sciences | Stellar astronomy | Astronomy |
536032 | https://en.wikipedia.org/wiki/Developmental%20disorder | Developmental disorder | Developmental disorders comprise a group of psychiatric conditions originating in childhood that involve serious impairment in different areas. There are several ways of using this term. The most narrow concept is used in the category "Specific Disorders of Psychological Development" in the ICD-10. These disorders comprise developmental language disorder, learning disorders, developmental coordination disorders, and autism spectrum disorders (ASD). In broader definitions, attention deficit hyperactivity disorder (ADHD) is included, and the term used is neurodevelopmental disorders. Yet others include antisocial behavior and schizophrenia that begins in childhood and continues through life. However, these two latter conditions are not as stable as the other developmental disorders, and there is not the same evidence of a shared genetic liability.
Developmental disorders are present from early life onward. Most improve as the child grows older, but some entail impairments that continue throughout life. These disorders differ from Pervasive developmental disorders (PPD), which uniquely describe a group of five developmental diagnoses, one of which is autism spectrum disorders (ASD). Pervasive developmental disorders reference a limited number of conditions whereas development disorders are a broad network of social, communicative, physical, genetic, intellectual, behavioral, and language concerns and diagnoses.
Emergence
Learning disabilities are often diagnosed when the children are young and just beginning school. Most learning disabilities are found under the age of 9.
Young children with communication disorders may not speak at all, or may have a limited vocabulary for their age. Some children with communication disorders have difficulty understanding simple directions or cannot name objects. Most children with communication disorders can speak by the time they enter school, however, they continue to have problems with communication. School-aged children often have problems understanding and formulating words. Teens may have more difficulty with understanding or expressing abstract ideas.
Causes
The scientific study of the causes of developmental disorders involves many theories. Some of the major differences between these theories involves whether environment disrupts normal development, if abnormalities are pre-determined, or if they are products of human evolutionary history which become disorders in modern environments (see evolutionary psychiatry).
Normal development occurs with a combination of contributions from both the environment and genetics. The theories vary in the part each factor has to play in normal development, thus affecting how the abnormalities are caused.
One theory that supports environmental causes of developmental disorders involves stress in early childhood. Researcher and child psychiatrist Bruce D. Perry, M.D., Ph.D, theorizes that developmental disorders can be caused by early childhood traumatization. In his works, he compares developmental disorders in traumatized children to adults with post-traumatic stress disorder, linking extreme environmental stress to the cause of developmental difficulties. Other stress theories suggest that even small stresses can accumulate to result in emotional, behavioral, or social disorders in children.
A 2017 study tested all 20,000 genes in about 4,300 families with children with rare developmental difficulties in the UK and Ireland in order to identify if these difficulties had a genetic cause. They found 14 new developmental disorders caused by spontaneous genetic mutations not found in either parent (such as a fault in the CDK13 gene). They estimated that about one in 300 children are born with spontaneous genetic mutations associated with rare developmental disorders.
Types
Autism spectrum disorder (ASD)
Diagnosis
The first diagnosed case of ASD was published in 1943 by American psychiatrist Leo Kanner. There is a wide range of cases and severity to ASD so it is very hard to detect the first signs of ASD. A diagnosis of ASD can be made accurately before the child is 3 years old, but the diagnosis of ASD is not commonly confirmed until the child is somewhat older. The age of diagnosis can range from 9 months to 14 years, and the mean age is 4 years old in the USA. On average each case of ASD is tested at three different diagnostic centers before confirmed. Early diagnosis of the disorder can diminish familial stress, speed up referral to special educational programs and influence family planning. The occurrence of ASD in one child can increase the risk of the next child having ASD by 50 to 100 times.
Abnormalities in the brain
The cause of ASD is still uncertain. What is known is that a child with ASD has a pervasive problem with how the brain is wired. Genes related to neurotransmitter receptors (serotonin and gamma-aminobutyric acid [GABA]) and CNS structural control (HOX genes) are found to be potential target genes that get affected in ASD. Autism spectrum disorder is a disorder of the many parts of the brain. Structural changes are observed in the cortex, which controls higher functions, sensation, muscle movements, and memory. Structural defects are seen in the cerebellum too, which affect the motor and communication skills. Sometimes the left lobe of the brain is affected and this causes neuropsychological symptoms. The distribution of white matter, the nerve fibers that link diverse parts of the brain, is abnormal. The corpus callosum, the band of nerve fibers, that connects the left and right hemispheres of the brain also gets affected in ASD. A study also found that 33% of people who have AgCC (agenesis of the corpus callosum), a condition in which the corpus callosum is partially or completely absent, had scores higher than the autism screening cut-off.
An ASD child's brain grows at a very rapid rate and is almost fully grown by the age of 10. Recent fMRI studies have also found altered connectivity within the social brain areas due to ASD and may be related to the social impairments encountered in ASD.
Symptoms
The symptoms have a wide range of severity. The symptoms of ASD can be broadly categorised as the following:
Persistent issues in social interactions and communications
These are predominantly seen by unresponsiveness in conversations, lesser emotional sharing, inability to initiate conversations, inability to interpret body language, avoidance of eye-contact and difficulty maintaining relationships.
Repetitive behavioral patterns
These patterns can be seen in the form of repeated movements of the hand or the phrases used while talking. A rigid adherence to schedules and inflexibility to adapt even if a minor change is made to their routine is also one of the behavioral symptoms of ASD. They could also display sensory patterns such as extreme aversion to certain odors or indifference to pain or temperature.
There are also different symptoms at different ages based on developmental milestones. Children between 0 and 36 months with ASD show a lack of eye contact, seem to be deaf, lack a social smile, do not like being touched or held, have unusual sensory behavior and show a lack of imitation. Children between 12 and 24 months with ASD show a lack of gestures, prefer to be alone, do not point to objects to indicate interest, are easily frustrated with challenges, and lack of functional play. And finally children between the ages 24 to 36 months with ASD show a lack of symbolic play and an unusual interest in certain objects, or moving objects.
Treatment
There is no specific treatment for autism spectrum disorders, but there are several types of therapy effective in easing the symptoms of autism, such as Applied Behavior Analysis (ABA), Speech-language therapy, Occupational therapy or Sensory integration therapy.
Applied behavioral analysis (ABA) is considered the most effective therapy for Autism spectrum disorders by the American Academy of Pediatrics. ABA focuses on teaching adaptive behaviors like social skills, play skills, or communication skills and diminishing problematic behaviors like self-injury. This is done by creating a specialized plan that uses behavioral therapy techniques, such as positive or negative reinforcement, to encourage or discourage certain behaviors over-time.
Occupational therapy helps autistic children and adults learn everyday skills that help them with daily tasks, such as personal hygiene and movement. These skills are then integrated into their home, school, and work environments. Therapists will oftentimes help patients learn to adapt their environment to their skill level. This type of therapy could help autistic people become more engaged in their environment. An occupational therapist will create a plan based on the patient's needs and desires and work with them to achieve their set goals.
Speech-language therapy can help those with autism who need to develop or improve communication skills. According to the organization Autism Speaks, “speech-language therapy is designed to coordinate the mechanics of speech with the meaning and social use of speech”. People with low-functioning autism may not be able to communicate with spoken words. Speech-language Pathologists (SLP) may teach someone how to communicate more effectively with others or work on starting to develop speech patterns. The SLP will create a plan that focuses on what the child needs.
Sensory integration therapy helps people with autism adapt to different kinds of sensory stimuli. Many children with autism can be oversensitive to certain stimuli, such as lights or sounds, causing them to overreact. Others may not react to certain stimuli, such as someone speaking to them. Many types of therapy activities involve a form of play, such as using swings, toys and trampolines to help engage the patients with sensory stimuli. Therapists will create a plan that focuses on the type of stimulation the person needs integration with.
Attention deficit hyperactivity disorder (ADHD)
Attention deficit hyperactivity disorder is a neurodevelopmental disorder that occurs in early childhood. ADHD affects 8 to 11% of children in the school going age. ADHD is characterised by significant levels of hyperactivity, inattentiveness, and impulsiveness. There are three subtypes of ADHD: predominantly inattentive, predominantly hyperactive, and combined (which presents as both hyperactive and inattentive subtypes). ADHD is twice as common in boys than girls but it is seen that the hyperactive/impulsive type is more common in boys while the inattentive type affects both sexes equally.
Symptoms
Symptoms of ADHD include inattentiveness, impulsiveness, and hyperactivity. Many of the behaviors that are associated with ADHD include poor control over actions resulting in disruptive behavior and academic problems. Another area that is affected by these disorders is the social arena for the person with the disorder. Many children that have this disorder exhibit poor interpersonal relationships and struggle to fit in socially with their peers. Behavioral study of these children can show a history of other symptoms such as temper tantrums, mood swings, sleep disturbances and aggressiveness.
Treatment options
The treatment of Attention Deficit Hyperactivity Disorder (ADHD) commonly involves a multimodal approach, combining various strategies to address the complex nature of the disorder. This comprehensive approach includes psychological, behavioral, pharmaceutical, and educational interventions tailored to the individual's specific needs. Here's a breakdown of the different components:
Psychological Interventions:
Counseling and Psychoeducation - Individuals with ADHD may benefit from counseling sessions that provide a safe space to discuss challenges, develop coping strategies, and improve self-esteem. **Psychoeducation helps individuals and their families understand the nature of ADHD and learn effective management techniques.
Cognitive Behavioral Therapy (CBT) - CBT aims to modify negative thought patterns and behaviors associated with ADHD. It helps individuals develop organizational skills, time management, and problem-solving abilities.
Behavioral Interventions:
Parent Training - Parents often participate in training programs to learn behavior management techniques. This may involve setting clear expectations, using positive reinforcement, and implementing consistent consequences for behavior.
Behavioral Modification Programs - These programs focus on shaping positive behaviors and reducing impulsive or disruptive behaviors in various settings, including home and school.
Pharmaceutical Interventions:
Stimulant Medications - Stimulant medications, such as methylphenidate (e.g., Ritalin) and amphetamines (e.g., Adderall), are commonly prescribed to manage symptoms of ADHD. These medications enhance the activity of neurotransmitters like dopamine and norepinephrine, helping to improve attention and impulse control.
Non-stimulant Medications - In cases where stimulants are not suitable or effective, non-stimulant medications like atomoxetine (Strattera) or guanfacine (Intuniv) may be prescribed.
Educational Interventions:
Individualized Education Plans (IEPs) - In educational settings, IEPs are developed to accommodate the unique learning needs of students with ADHD. This may involve classroom modifications, additional support, and specific teaching strategies.
504 Plans - These plans outline accommodations for students with ADHD in mainstream educational settings, such as extended test-taking time or preferential seating.
The effectiveness of the treatment plan depends on the individual's specific challenges and responses to interventions. A collaborative and multidisciplinary approach involving parents, educators, mental health professionals, and healthcare providers is crucial for developing and implementing a successful ADHD management plan. Regular monitoring and adjustments to the treatment plan may be necessary to meet the evolving needs of individuals with ADHD.
Behavioral therapy
Sessions of counselling, cognitive behavioral therapy (CBT), making environmental changes in noise and visual stimulation are some behavioral management techniques followed. But it has been observed that behavioral therapy alone is less effective than therapy with stimulant drugs alone.
Drug therapy
Medications commonly utilized in the treatment of Attention Deficit Hyperactivity Disorder (ADHD) include stimulants like methylphenidate and lisdexamfetamine, as well as non-stimulants such as atomoxetine. These medications can effectively manage ADHD symptoms by targeting neurotransmitter imbalances. However, it is important to be aware of potential side effects associated with these medications. Common side effects may include headaches, which can often be mitigated by adjusting the dosage or administration timing. Gastrointestinal discomfort, including stomach pain or nausea, is another possible side effect, and taking the medication with food or modifying the dosage may help alleviate these symptoms. Additionally, while rare, changes in mood such as feelings of depression have been reported. Careful monitoring and communication with healthcare providers are essential to address and manage any side effects, ensuring the overall effectiveness and well-being of individuals undergoing ADHD treatments.
SSRI antidepressants may be unhelpful, and could worsen symptoms of ADHD. However ADHD is often misdiagnosed as depression, particularly when no hyperactivity is present.
Other disorders
Learning disabilities, such as Dysgraphia
Communication disorders and Auditory processing disorder
Developmental coordination disorder
Genetic disorders, such as Down syndrome or Williams syndrome
Tic disorders such as Tourette syndrome
Stuttering
Intellectual disability
| Biology and health sciences | Mental disorders | Health |
536059 | https://en.wikipedia.org/wiki/Qanat | Qanat | A qanāt () or kārīz () is a water supply system that was developed in ancient Iran for the purpose of transporting usable water to the surface from an aquifer or a well through an underground aqueduct. Originating approximately 3,000 years ago, its function is essentially the same across the Middle East and North Africa, but it is known by a variety of regional names beyond today's Iran, including: kārēz in Afghanistan and Pakistan; foggāra in Algeria; khettāra in Morocco; falaj in Oman and the United Arab Emirates; and ʿuyūn in Saudi Arabia. In addition to those in Iran, the largest extant and functional qanats are located in Afghanistan, Algeria, China (i.e., the Turpan water system), Oman, and Pakistan.
Proving crucial to water supply in areas with hot and dry climates, a qanat enables water to be transported over long distances by largely eliminating the risk of much of it evaporating on the journey. The system also has the advantage of being fairly resistant to natural disasters, such as floods and earthquakes, as well as to man-made disasters, such as wartime destruction and water supply terrorism. Furthermore, it is almost insensitive to varying levels of precipitation, delivering a flow with only gradual variations from wet to dry years.
The typical design of a qanat is that of a series of well-like vertical shafts, which are all connected by a gently sloping tunnel. This taps into groundwater and delivers it to the surface via gravity, therefore eliminating the need for pumping. The vertical shafts along the underground channel are for maintenance purposes, and water is typically used only once it emerges from the daylight point.
To date, the qanat system still ensures a reliable supply of water for consumption and irrigation across human settlements in hot, arid, and semi-arid climates, but its value to a population is directly related to the quality, volume, and regularity of the groundwater in the inhabited region. Since their adoption outside of the Iranian mainland in antiquity, qanats have come to be heavily relied upon by much of the Middle Eastern and North African populations for sustenance. Likewise, many of the continuously inhabited settlements in these regions are established in areas where conditions have historically been favourable for creating and sustaining a qanat system.
Names
Common variants of qanat in English include kanat, khanat, kunut, kona, konait, ghanat, ghundat.
() is an Arabic word that means "channel". In Persian, the words for "qanat" are (or ; ) and is derived from earlier word (). The word () is also used in Persian. Other names for qanat include (); (Balochi); (Azerbaijan); (Morocco); , or (Spain); () (United Arab Emirates and Oman); (North Africa). Alternative terms for qanats in Asia and North Africa are kakuriz, chin-avulz, and mayun.
Origins
According to most sources, qanat technology was developed by the ancient Iranians sometime in the early 1st millennium BCE and slowly spread westward and eastward from there. Other sources suggest a Southeast Arabian origin. Analogous systems appear to have been developed independently in China and in South America (specifically, southern Peru).
A cotton species, Gossypium arboreum, is indigenous to South Asia and has been cultivated on the Indian subcontinent for a long time. Cotton appears in the Inquiry into Plants by Theophrastus and is mentioned in the Laws of Manu. As transregional trade networks expanded and intensified, cotton spread from its homeland to India and into the Middle East. One theory is that the qanat was developed to irrigate cotton fields, first in what is now Iran, where it doubled the amount of available water for irrigation and urban use. Because of this, Persia enjoyed larger surpluses of agricultural products, thus increasing urbanization and social stratification. The qanat technology subsequently spread from Persia westward and eastward.
In the arid coastal desert of Peru, a technology of water supply similar to that of the qanats, called puquios, was developed. Most archaeologists believe that the puquios are indigenous and date to about 500 CE, but a few believe they are of Spanish origin, brought to the Americas in the 16th century. Puquios were still in use in the Nazca region in the 21st century.
Features
Qanats are constructed as a series of well-like vertical shafts, connected by a gently sloping tunnel which carries a water canal. Qanats efficiently deliver large amounts of subterranean water to the surface without need for pumping. The water drains by gravity, typically from an upland aquifer, with the destination lower than the source. Qanats allow water to be transported over long distances in hot dry climates without much water loss to evaporation.
It is very common for a qanat to start below the foothills of mountains, where the water table is closest to the surface. From this source, the qanat tunnel slopes gently downward, slowly converging with the steeper slope of the land surface above, and the water finally flows out above ground where the two levels meet. To connect a populated or agricultural area with an aquifer, qanats must often extend for long distances.
Qanats are sometimes split into an underground distribution network of smaller canals called kariz. Like qanats, these smaller canals are below ground to avoid contamination and evaporation. In some cases water from a qanat is stored in a reservoir, typically with night flow stored for daytime use. An ab anbar is an example of a traditional Persian qanat-fed reservoir for drinking water.
The qanat system has the advantage of being resistant to natural disasters such as floods, and to deliberate destruction in war. Furthermore, it is almost insensitive to the levels of precipitation, delivering a flow with only gradual variations from wet to dry years. From a sustainability perspective, qanats are powered only by gravity and thus have low operation and maintenance costs. Qanats transfer fresh water from the mountain plateau to the lower-lying plains with saltier soil. This helps to control soil salinity and prevent desertification.
The qanat should not be confused with the spring-flow tunnel typical to the mountainous area around Jerusalem. Although both are excavated tunnels designed to extract water by gravity flow, there are crucial differences. Firstly, the origin of the qanat was a well that was turned into an artificial spring. In contrast, the origin of the spring-flow tunnel was the development of a natural spring to renew or increase flow following a recession of the water table. Secondly, the shafts essential for the construction of qanats are not essential to spring-flow tunnels.
Impact on settlement patterns
A typical town or city in Iran, and elsewhere where the qanat is used, has more than one qanat. Fields and gardens are located both over the qanats a short distance before they emerge from the ground and below the surface outlet. Water from the qanats define both the social regions in the city and the layout of the city.
The water is freshest, cleanest, and coolest in the upper reaches, and more prosperous people live at the outlet or immediately upstream of the outlet. When the qanat is still below ground, the water is drawn to the surface via wells or animal driven Persian wells. Private subterranean reservoirs could supply houses and buildings for domestic use and garden irrigation as well. Air flow from the qanat is used to cool an underground summer room (shabestan) found in many older houses and buildings.
Downstream of the outlet, the water runs through surface canals called jubs (jūbs) which run downhill, with lateral branches to carry water to the neighborhood, gardens and fields. The streets normally parallel the jubs and their lateral branches. As a result, the cities and towns are oriented consistent with the gradient of the land; this is a practical response to efficient water distribution over varying terrain.
The lower reaches of the canals are less desirable for both residences and agriculture. The water grows progressively more polluted as it passes downstream. In dry years the lower reaches are the most likely to see substantial reductions in flow.
Construction
Traditionally qanats are built by a group of skilled laborers, muqannīs, with hand labor. The profession historically paid well and was typically handed down from father to son.
Preparations
The critical, initial step in qanat construction is identification of an appropriate water source. The search begins at the point where the alluvial fan meets the mountains or foothills; water is more abundant in the mountains because of orographic lifting, and excavation in the alluvial fan is relatively easy. The muqannīs follow the track of the main water courses coming from the mountains or foothills to identify evidence of subsurface water such as deep-rooted vegetation or seasonal seeps. A trial well is then dug to determine the depth of the water table and determine whether a sufficient flow is available to justify construction. If these prerequisites are met, the route is laid out aboveground.
Equipment must be assembled. The equipment is straightforward: containers (usually leather bags), ropes, reels to raise the container to the surface at the shaft head, hatchets and shovels for excavation, lights, and spirit levels or plumb bobs and string. Depending upon the soil type, qanat liners (usually fired clay hoops) may also be required.
Although the construction methods are simple, the construction of a qanat requires a detailed understanding of subterranean geology and a degree of engineering sophistication. The gradient of the qanat must be carefully controlled: too shallow a gradient yields no flow and too steep a gradient will result in excessive erosion, collapsing the qanat. And misreading the soil conditions leads to collapses, which at best require extensive rework and at worst are fatal for the crew.
Excavation
Construction of a qanat is usually performed by a crew of 3–4 muqannīs. For a shallow qanat, one worker typically digs the horizontal shaft, one raises the excavated earth from the shaft and one distributes the excavated earth at the top.
The crew typically begins from the destination to which the water will be delivered into the soil and works toward the source (the test well). Vertical shafts are excavated along the route, separated at a distance of . The separation of the shafts is a balance between the amount of work required to excavate them and the amount of effort required to excavate the space between them, as well as the ultimate maintenance effort. In general, the shallower the qanat, the closer the vertical shafts. If the qanat is long, excavation may begin from both ends at once. Tributary channels are sometimes also constructed to supplement the water flow.
Most qanats in Iran run less than , while some have been measured at ≈ in length near Kerman. The vertical shafts usually range from in depth, although qanats in the province of Khorasan have been recorded with vertical shafts of up to . The vertical shafts support construction and maintenance of the underground channel as well as air interchange. Deep shafts require intermediate platforms to facilitate the process of removing soil.
The construction speed depends on the depth and nature of the ground. If the earth is soft and easy to work, at depth a crew of four workers can excavate a horizontal length of per day. When the vertical shaft reaches , they can excavate only 20 meters horizontally per day and at in depth this drops below 5 horizontal meters per day. In Algeria, a common speed is just per day at a depth of . Deep, long qanats (which many are) require years and even decades to construct.
The excavated material is usually transported by means of leather bags up the vertical shafts. It is mounded around the vertical shaft exit, providing a barrier that prevents windblown or rain driven debris from entering the shafts. These mounds may be covered to provide further protection to the qanat. From the air, these shafts look like a string of bomb craters.
The qanat's water-carrying channel must have a sufficient downward slope that water flows easily. However the downward gradient must not be so great as to create conditions under which the water transitions between supercritical and subcritical flow. If this occurs, the waves that result can result in severe erosion that can damage or destroy the qanat. The choice of the slope is a trade off between erosion and sedimentation. Highly sloped tunnels are subject to more erosion as water flows at a higher speed. On the other hand, less sloped tunnels need frequent maintenance due to the problem of sedimentation. A lower downward gradient also contributes to reducing the solid contents and contamination in water. In shorter qanats the downward gradient varies between 1:1000 and 1:1500, while in longer qanats it may be almost horizontal. Such precision is routinely obtained with a spirit level and string.
In cases where the gradient is steeper, underground waterfalls may be constructed with appropriate design features (usually linings) to absorb the energy with minimal erosion. In some cases the water power has been harnessed to drive underground mills. If it is not possible to bring the outlet of the qanat out near the settlement, it is necessary to run a jub or canal overground. This is avoided when possible to limit pollution, warming and water loss due to evaporation.
Maintenance
The vertical shafts may be covered to minimize blown-in sand. The channels of qanats must be periodically inspected for erosion or cave-ins, cleaned of sand and mud and otherwise repaired. For safety, air flow must be assured before entry.
Some damaged qanats have been restored. To be sustainable, restoration needs to take into account many nontechnical factors beginning with the process of selecting the qanat to be restored. In Syria, three sites were chosen based on a national inventory conducted in 2001. One of them, the Drasiah qanat of Dmeir, was completed in 2002. Selection criteria included the availability of a steady groundwater flow, social cohesion and willingness to contribute of the community using the qanat, and the existence of a functioning water-rights system.
Applications
The primary applications of qanats are for irrigation, providing cattle with water, and drinking water supply. Other applications include watermills, cooling and ice storage.
Watermills
Watermills within a qanat system had to be carefully situated, to make best use of the slow flow of water. In Iran, there were subterranean mills at Yazd and Boshruyeh; at Taft and Ardestan mills were placed at the outflow from the qanat, before irrigation of the fields.
Cooling
Qanats used in conjunction with a wind tower can provide cooling as well as a water supply. A wind tower is a chimney-like structure positioned above the house; of its four openings, the one opposite the wind direction is opened to move air out of the house. Incoming air is pulled from a qanat below the house. The air flow across the vertical shaft opening creates a lower pressure (see Bernoulli effect) and draws cool air up from the qanat tunnel, mixing with it.
The air from the qanat is drawn into the tunnel at some distance away and is cooled both by contact with the cool tunnel walls/water and by the transfer of latent heat of evaporation as water evaporates into the air stream. In dry desert climates this can result in a greater than 15 °C reduction in the air temperature coming from the qanat; the mixed air still feels dry, so the basement is cool and only comfortably moist (not damp). Wind tower and qanat cooling have been used in desert climates for over 1,000 years.
Ice storage
By 400 BCE, Persian engineers had mastered the technique of storing ice in the middle of summer in the desert.
The ice could be brought in during the winters from nearby mountains, but in a more usual and sophisticated method they built a wall in the east–west direction near a yakhchal (ice pit). In winter, the qanat water would be channeled to the north side of the wall, whose shade made the water freeze more quickly, increasing the ice formed per winter day. Then the ice was stored in yakhchals—specially designed, naturally cooled refrigerators.
A large underground space with thick insulated walls was connected to a qanat, and a system of windcatchers or wind towers was used to draw cool subterranean air up from the qanat to maintain temperatures inside the space at low levels, even during hot summer days. As a result, the ice melted slowly and was available year-round.
By country
Africa
Algeria
Qanats (designated foggaras in Algeria) are the source of water for irrigation in large oases like Gourara. The foggaras are also found at Touat (an area of Adrar 200 km from Gourara). The length of the foggaras in this region is estimated to be thousands of kilometers. Although sources suggest that the foggaras may have been in use as early as 200 CE, they were clearly in use by the 11th century after the Arabs took possession of the oases in the 10th century and the residents embraced Islam. The water is metered to the various users through the use of distribution weirs that meter flow to the various canals, each for a separate user.
The humidity of the oases is also used to supplement the water supply to the foggara. The temperature gradient in the vertical shafts causes air to rise by natural convection, causing a draft to enter the foggara. The moist air of the agricultural area is drawn into the foggara in the opposite direction to the water run-off. In the foggara it condenses on the tunnel walls and the air passes out of the vertical shafts. This condensed moisture is available for reuse.
Egypt
Qanat irrigation technology was introduced to Egypt by the Achaemenid king Darius I during his reign of 522 BCE-486 BCE, which is supported by the historian Albert T. Olmstead. There are four main oases in the Egyptian desert. The Kharga Oasis is one that has been extensively studied. There is evidence that as early as the second half of the 5th century BCE water brought in qanats was being used. The qanats were excavated through water-bearing sandstone rock, which seeps into the channel, with water collected in a basin behind a small dam at the end. The width is approximately , but the height ranges from 5 to 9 meters; it is likely that the qanat was deepened to enhance seepage when the water table dropped (as is also seen in Iran). From there the water was used to irrigate fields.
There is another instructive structure located at the Kharga Oasis. A well that apparently dried up was improved by driving a side shaft through the easily penetrated sandstone (presumably in the direction of greatest water seepage) into the hill of Ayn-Manâwîr (also written to allow collection of additional water. After this side shaft had been extended, another vertical shaft was driven to intersect the side shaft. Side chambers were built, and holes bored into the rock—presumably at points where water seeped from the rocks—are evident.
Libya
David Mattingly reports foggara extending for hundreds of miles in the Garamantes area near Germa in Libya: "The channels were generally very narrow – less than 2 feet wide and 5 high – but some were several miles long, and in total some 600 foggara extended for hundreds of miles underground. The channels were dug out and maintained using a series of regularly spaced vertical shafts, one every 30 feet or so, 100,000 in total, averaging 30 feet in depth, but sometimes reaching 130."
Morocco
In southern Morocco, the qanat (locally khettara) is also used. On the margins of the Sahara Desert, the isolated oases of the Draa River valley and Tafilalt have relied on qanat water for irrigation since the late 14th century. In Marrakech and the Haouz plain, the qanats have been abandoned since the early 1970s, having dried up. In the Tafilaft area, half of the 400 khettaras are still in use. The 1971 Hassan Adahkil Dam's build in the main course of the Ziz River and its subsequent impact on local water tables is said to be one of the many reasons for the loss of half of the khettara.
The black berbers (haratin) of the south were the hereditary class of qanat diggers in Morocco who build and repair these systems. Their work was hazardous.
Tunisia
The foggara water management system in Tunisia, used to create oases, is similar to that of the Iranian qanat. The foggara is dug into the foothills of a fairly steep mountain range such as the eastern ranges of the Atlas Mountains. Rainfall in the mountains enters the aquifer and moves toward the Saharan region to the south. The foggara, in length, penetrates the aquifer and collects water. Families maintain the foggara and own the land it irrigates over a ten-meter width, with length reckoned by the size of plot that the available water will irrigate.
Asia
Afghanistan
The qanats are called kariz in Dari (Persian) and Pashto and have been in use since the pre-Islamic period. It is estimated that more than 9,370 karizes were in use in the 20th century. The oldest functional kariz which is more than 300 years old and 8 kilometers long is located in Wardak province and is still providing water to nearly 3,000 people.
Many of these ancient structures were destroyed during the Soviet–Aghan War and the War in Afghanistan. Maintenance has not always been possible. The cost of labour has become very high, and maintaining the kariz structures is no longer possible. Lack of skilled artisans who have the traditional knowledge also poses difficulties. A number of the large farmers are abandoning their kariz which has been in their families sometimes for centuries, and moving to tube and dug wells backed by diesel pumps. However, the government of Afghanistan was aware of the importance of these structures and all efforts were made to repair, reconstruct and maintain (through the community) the kariz. The Ministry of Rural Rehabilitation and Development along with national and international NGOs made the effort. There were still functional qanat systems in 2009. American forces were reported to have unintentionally destroyed some of the channels during expansion of a military base, creating tensions between them and the local community. Some of these tunnels were used to store supplies, and to move men and equipment underground.
Armenia
Qanats have been preserved in Armenia in the community of Shvanidzor, in the southern province of Syunik, bordering with Iran. Qanats are named kahrezes in Armenian. There are 5 kahrezes in Shvanidzor. Four of them were constructed before the village was founded. The fifth kahrez was constructed in 2005. Potable water runs through three of them, and two are in poor condition. In the summer, especially in July and August, the amount of water reaches its minimum, creating a critical situation in the water supply system. Still, kahrezes are the main source of potable and irrigation water for the community.
Azerbaijan
The territory of Azerbaijan was home to numerous kahrizes many centuries ago. Archaeological findings suggest that long before the 9th century CE, kahrizes by which the inhabitants brought potable and irrigation water to their settlements were in use in Azerbaijan. Traditionally, kahrizes were built and maintained by a group of masons called 'Kankans' with manual labour. The profession was handed down from father to son.
It is estimated that until the 20th century, nearly 1,500 kahrizes, of which as many as 400 were in the Nakhichevan Autonomous Republic, existed in Azerbaijan. However, following the introduction of electric and fuel-pumped wells during Soviet times, kahrizes were neglected. Today, it is estimated that 800 are still functioning in Azerbaijan. These operational kahrizes are key to the life of many communities.
In 1999, upon the request of the communities in Nakhichevan, the International Organization for Migration (IOM) began implementing a pilot programme to rehabilitate the kahrizes. By 2018 IOM rehabilitated more than 163 kahrizes with funds from the United Nations Development Programme, European Commission, Canadian International Development Agency, Swiss Agency for Development and Cooperation and the Bureau of Population, Refugees, and Migration, US State Department, and the self-contribution of the local communities.
In 2010, IOM began a kahriz rehabilitation project with funds from the Korea International Cooperation Agency. During the First Phase of the action which lasted until January 2013, a total of 20 kahrizes in the mainland of Azerbaijan have been renovated.
China
The oasis of Turpan, in the deserts of Xinjiang in northwestern China, uses water provided by qanat (locally called karez). There are nearly 1,000 karez systems in the area, and the total length of the canals is about 5,000 kilometers.
Turpan has long been the center of a fertile oasis and an important trade center along the Northern Silk Road, at which time it was adjacent to the kingdoms of Korla and Karashahr to the southwest. The historical record of the karez extends back to the Han dynasty. The Turfan Water Museum is a Protected Area of the People's Republic of China because of the importance of the Turpan karez water system to the history of the area.
Iran
In the middle of the 20th century, an estimated 50,000 qanats were in use in Iran, each commissioned and maintained by local users. Of these, only 37,000 remain in use as of 2015. One of the oldest and largest known qanats is in the Iranian city of Gonabad, and after 2,700 years still provides drinking and agricultural water to nearly 40,000 people. Its main well depth is more than 360 meters and its length is 45 kilometers. Yazd, Khorasan and Kerman are zones known for their dependence on an extensive system of qanats.
In 2016, UNESCO inscribed the Persian Qanat as a World Heritage Site, listing the following eleven qanats: Qasebeh Qanat, Qanat of Baladeh, Qanat of Zarch, Hasan Abad-e Moshir Qanat, Ebrāhim Ābād Qanat in Markazi Province, Qanat of Vazvān in Esfahan Province, Mozd Ābād Qanat in Esfahan Province, Qanat of the Moon in Esfahan Province, Qanat of Gowhar-riz in Kerman Province, Jupār – Ghāsem Ābād Qanat in Kerman Province, and Akbar Ābād Qanat in Kerman Province. Since 2002, UNESCO's International Hydrological Programme Intergovernmental Council began investigating the possibility of an international qanat research center to be located in Yazd, Iran.
The Qanats of Gonabad, also called kariz Kai Khosrow, is one of the oldest and largest qanats in the world built between 700 BCE to 500 BCE. It is located at Gonabad, Razavi Khorasan Province. This property contains 427 water wells with total length of .
According to Callisthenes, the Persians were using water clocks in 328 BCE to ensure a just and exact distribution of water from qanats to their shareholders for agricultural irrigation. The use of water clocks in Iran, especially in Qanats of Gonabad and kariz Zibad, dates back to 500 BCE. Later they were also used to determine the exact holy days of pre-Islamic religions, such as the Nowruz, Chelah, or Yaldā – the shortest, longest, and equal-length days and nights of the years.
The water clock, or Fenjaan, was the most accurate and commonly used timekeeping device for calculating the amount or the time that a farmer must take water from the Qanats of Gonabad until it was replaced by more accurate current clocks.
Many of the Iranian qanats bear some characteristics which allow us to call them feat of engineering, considering the intricate techniques used in their construction. The eastern and central regions of Iran hold the most qanats due to low precipitation and lack of permanent surface streams, whereas a small number of qanats can be found in the northern and western parts which receive more rainfall and enjoy some permanent rivers. Respectively the provinces Khorasan Razavi, Southern Khorasan, Isfahan, and Yazd accommodate the most qanats, but from the viewpoint of water discharge the provinces Isfahan, Khorasan Razavi, Fars and Kerman are ranked first to fourth.
Henri Golbot explored the genesis of the qanat in his 1979 publication, (The Qanats. A Technique for Obtaining Water), He argues that the ancient Iranians made use of the water that the miners wished to get rid of it, and founded a basic system named qanat or kariz to supply the required water to their farm lands. According to Golbot, this innovation took place in the northwest of the present Iran somewhere bordering Turkey and later was introduced to the neighboring Zagros Mountains.
According to an inscription left by Sargon II, the king of Assyria, in 714 BCE he invaded the city of Uhlu lying in the northwest of Uroomiye lake that lay in the territory of Urartu empire, and then he noticed that the occupied area enjoyed a very rich vegetation even though there was no river running across it. So he managed to discover the reason why the area could stay green and realized that there were some qanats behind the matter. In fact it was Ursa, the king of the region, who had rescued the people from thirst and turned Uhlu into a prosperous and green land. Golbot believes that the influence of the Medeans and Achaemenids made the technology of qanat spread from Urartu, in the western north of Iran and near the present border between Iran and Turkey, to all over the Iranian plateau.
It was an Achaemenid ruling that in case someone succeeded in constructing a qanat and bringing groundwater to the surface in order to cultivate land, or in renovating an abandoned qanat, the tax he was supposed to pay the government would be waived not only for him but also for his successors for up to 5 generations. During this period, the technology of qanat was in its heyday and it even spread to other countries. For example, following Darius's order, Silaks the naval commander of the Persian army and Khenombiz the royal architect managed to construct a qanat in the oasis of Kharagha in Egypt.
Beadnell believes that qanat construction dates back to two distinct periods: they were first constructed by the Persians, and later the Romans dug some other qanats during their reign in Egypt from 30 BCE to 395 CE. The magnificent temple built in this area during Darius's reign shows that there was a considerable population depending on the water of qanats. Ragerz has estimated this population to be 10,000 people. The most reliable document confirming the existence of qanats at this time was written by Polybius who states that: "the streams are running down from everywhere at the base of Alborz mountain, and people have transferred too much water from a long distance through some subterranean canals by spending much cost and labor."
During the Seleucid era, which began after the occupation of Iran by Alexander the Great, it seems that the qanats were abandoned. In terms of the situation of qanats during this era, some historical records have been found. In a study by Russian orientalist scholars it has been mentioned that: the Persians used the side branches of rivers, mountain springs, wells and qanats to supply water. The subterranean galleries excavated to obtain groundwater were named as qanat. These galleries were linked to the surface through some vertical shafts which were sunk in order to get access to the gallery to repair it if necessary.
According to the historical records, the Parthian kings did not care about the qanats the way the Achaemenid kings and even Sassanid kings did. As an instance, Arsac III, one of the Parthian kings, destroyed some qanats in order to make it difficult for Seleucid Antiochus to advance further while fighting him. The historical records from this time indicate a perfect regulation on both water distribution and farmlands. All the water rights were recorded in a special document which was referred to in case of any transaction. The lists of farmlands – whether private or governmental – were kept at the tax department. During this period there were some official rulings on qanats, streams, construction of dam, operation and maintenance of qanats, etc.
The government proceeded to repair or dredge the qanats that were abandoned or destroyed, and to construct the new qanats if necessary. A document written in the Pahlavi language points out the important role of qanats in developing the cities at that time. In Iran, the advent of Islam, which coincided with the overthrow of the Sassanid dynasty, brought about a profound change in religious, political, social and cultural structures. But the qanats stayed intact because the economic infrastructure including qanats was of great importance to the Arabs. As an instance, M. Lombard reports that the Moslem clerics who lived during Abbasid period, such as Abooyoosef Ya'qoob (died 798 CE) stipulated that whoever can bring water to the idle lands in order to cultivate, his tax would be waived and he would be entitled to the lands cultivated. Therefore, this policy did not differ from that of the Achaemenids in not getting any tax from the people who revived abandoned lands.
The Arabs' supportive policy on qanats was so successful that even Mecca gained a qanat. The Persian historian Hamdollah Mostowfi writes: "Zobeyde Khatoon (Haroon al-Rashid's wife) constructed a qanat in Mecca. After the time of Haroon al-Rashid, during the caliph Moghtader's reign this qanat fell into decay, but he rehabilitated it, and the qanat was rehabilitated again after it collapsed during the reign of two other caliphs named Ghaem and Naser. After the era of the caliphs this qanat completely fell into ruin because the desert sand filled it up, but later Amir Choopan repaired the qanat and made it flow again in Mecca."
There are also other historical texts proving that the Abbasids were concerned about qanats. For example, according to the "Incidents of Abdollah bin Tahir's Time" written by Gardizi, in 830 CE a terrible earthquake struck the town of Forghaneh and reduced many homes to rubble. The inhabitants of Neyshaboor used to come to Abdollah bin Tahir in order to request him to intervene, for they fought over their qanats and found the relevant instruction or law on qanat as a solution neither in the prophet's quotations nor in the clerics' writings. So Abdollah bin Tahir managed to bring together all the clergymen from throughout Khorasan and Iraq to compile a book entitled Alghani (The Book of Qanat). This book collected all the rulings on qanats which could be of use to whoever wanted to judge a dispute over this issue. Gardizi added that this book was still applicable to his time, and everyone made references to this book. One can deduce from these facts that during the above-mentioned period the number of qanats was so considerable that the authorities were prompted to put together some legal instructions concerning them. Also it shows that from the 9th to 11th centuries the qanats that were the hub of the agricultural systems were also of interest to the government.
Apart from "The Book of Alghani", which is considered as a law booklet focusing on qanat-related rulings based on Islamic principles, there is another book about groundwater written by Karaji in 1010. This book, entitled Extraction of Hidden Waters, examines just the technical issues associated with the qanat and tries to answer the common questions such as how to construct and repair a qanat, how to find a groundwater supply, how to do leveling, etc.. Some of the innovations described in this book were introduced for the first time in the history of hydrogeology, and some of its technical methods are still valid and can be applied in qanat construction. The content of the book implies that its writer (Karaji) did not have any idea that there was another book on qanats compiled by the clergymen.
There are some records dating back to that time, signifying their concern about the legal vicinity of qanats. For example, Mohammad bin Hasan quotes Aboo-Hanifeh that in case someone constructs a qanat in abandoned land, someone else can dig another qanat in the same land on the condition that the second qanat is 500 zera' (375 meters) away from the first one.
Ms. Lambton quotes Moeen al-din Esfarzi who wrote the book Rowzat al-Jannat (the garden of paradise) that Abdollah bin Tahir (from the Taherian dynasty) and Ismaeel Ahmed Samani (from the Samani dynasty) had several qanats constructed in Neyshaboor. Later, in the 11th century, a writer named Nasir Khosrow acknowledged all those qanats with the following words: "Neyshaboor is located in a vast plain at a distance of 40 Farsang (≈240 km) from Serakhs and 70 Farsang (≈420 km) from Mary (Marv) ... all the qanats of this city run underground, and it is said that an Arab who was offended by the people of Neyshaboor has complained that; what a beautiful city Neyshaboor could have become if its qanats would have flowed on the ground surface and instead its people would have been underground." These documents all certify the importance of qanats during the Islamic history within the cultural territories of Iran.
In the 13th century, the invasion of Iran by Mongolian tribes reduced many qanats and irrigation systems to ruin, and many qanats were deserted and dried up. Later, in the era of the Ilkhanid dynasty especially at the time of Ghazan Khan and his Persian minister Rashid al-Din Fazl-Allah, some measures were taken to revive the qanats and irrigation systems. There is a 14th-century book entitled Al-Vaghfiya Al-Rashidiya (Rashid's Deeds of Endowment) that names all the properties located in Yazd, Shiraz, Maraghe, Tabriz, Isfahan and Mowsel that Rashid Fazl-Allah donated to the public or religious places. This book mentions many qanats running at that time and irrigating a considerable area of farmland.
At the same time, another book, entitled Jame' al-Kheyrat, was written by Seyyed Rokn al-Din on the same subject as Rashid's book. In this book, Seyyed Rokn al-Din names the properties he donated in the region of Yazd. These deeds of endowment indicate that much attention was given to the qanats during the reign of Ilkhanids, but it is attributable to their Persian ministers, who influenced them.
In 1984–1985 the ministry of energy took a census of 28,038 qanats whose total discharge was 9 billion cubic meters. In the years 1992–1993 the census of 28,054 qanats showed a total discharge of 10 billion cubic meters. 10 years later in 2002–2003 the number of the qanats was reported as 33,691 with a total discharge of 8 billion cubic meters.
In the restricted regions there are 317,225 wells, qanats and springs that discharge 36,719 million cubic meters water per year, out of which 3,409 million cubic meters is surplus to the aquifer capacity. in 2005, in the country as a whole, there were 130,008 deep wells with a discharge of 31,403 million cubic meters, 33,8041 semi deep wells with a discharge of 13,491 million cubic meters, 34,355 qanats with a discharge of 8,212 million cubic meters, and 55,912 natural springs with a discharge of 21,240 million cubic meters.
In 2021, a British-trained architect Margot Krasojević designed a luxury eco hotel based on principles of qanat and windcatchers in a desert in Iran, called Qanat. The project has not yet been built but offers ideas for applying ancient technology to modern-day cooling problems in the desert.
Iraq
A survey of qanat systems in the Kurdistan region of Iraq conducted by the Department of Geography at Oklahoma State University (US) on behalf of UNESCO in 2009 found that out of 683 karez systems, some 380 were still active in 2004, but only 116 were active by 2009. Reasons for the decline of qanats include "abandonment and neglect" prior to 2004, "excessive pumping from wells" and, since 2005, drought. Water shortages are said to have forced, since 2005, over 100,000 people who depended for their livelihoods on karez systems to leave their homes.
The study says that a single karez has the potential to provide enough household water for nearly 9,000 individuals and irrigate over 200 hectares of farmland. UNESCO and the government of Iraq plan to rehabilitate the karez through a Karez Initiative for Community Revitalization launched in 2010. Most of the karez are in Sulaymaniyah Governorate (84%). A large number are also found in Erbil Governorate (13%), especially on the broad plain around and in Erbil city.
India
In India, there are karez systems are located at Bidar, Bijapur, Burhanpur "(Kundi Bhandara)", and Aurgangabad. The Bidar karez systems were probably the first dug in India. It dates to the Bahmani period. Bidar has three karez systems as per Ghulam Yazdani's documentation. Other than Naubad there are two more karez systems in Bidar, "Shukla Theerth" and "Jamna Mori". The Shukla theerth is the longest karez system in Bidar. The mother well of this karez has been discovered by near Gornalli Kere, a historic embankment. The third system called Jamna mori is more of a distribution system within the old city area with many channels crisscrossing the city lanes. Restoration efforts commenced in 2014, with the desilting and excavation of the Naubad Karez in 2015, uncovering 27 vertical shafts linked to the Karez. The rejuvenation of the system has had a significant impact on the water-deficit city of Bidar. A seventh line of the system was discovered in 2016 during a sewage line excavation.
Valliyil Govindankutty, assistant professor in geography at Government College, Chittur, was responsible for rediscovery and mapping of the Naubad Karez System in 2012-2013. Later in 2014-2016 team YUVAA joined Govindankutty to help uncover Other two Karez Systems in Bidar. Detailed documentation of the Naubad karez system was done in August 2013 and a report was submitted to District Administration of Bidar that found several new facts. The research has led to the initiation of cleaning the debris and collapsed sections paving the way to its rejuvenation. The cleaning of karez has led to bringing water to higher areas of the plateau, and it has in turn recharged the wells in the vicinity.
The Bijapur karez system is much more complicated. A reveals that it has surface water and groundwater connections. The Bijapur karez is a network of shallow masonry aqueducts, terracotta/ceramic pipes, embankments and reservoirs, tanks etc. All weave together a network to ensure water reaches the old city. The system starts at Torwi and extends as shallow aqueducts and further as pipes; further it becomes deeper from the Sainik school area onward which exists as a tunnel dug through the geology. The system can be clearly traced up to Ibrahim Roja.
In Aurangabad the karez systems are called nahars. These are shallow aqueducts running through the city. There are 14 aqueducts in Aurangabad. The Nahar-i-Ambari is the oldest and longest. Its again a combination of shallow aqueducts, open channels, pipes, cisterns, etc. The source of water is a surface water body. The karez has been constructed right below the bed of lake. The lake water seeps through the soil into the Karez Gallery.
In Burhanpur the karez is called "Kundi-Bhandara", sometimes wrongly referred to as"Khuni Bhandara". The system is approx 6 km long starts from the alluvial fans of Satpura hills in the north of the town. Unlike Bidar, Bijapur and Aurgangabad the System airvents are round in shape. Inside the Karez one could see lime depositions on the walls. The Systems ends to carry water further to palaces and public fountains through pipe line.
Indonesia
It has been suggested that underground temples at Gua Made in Java reached by shafts, in which masks of a green metal were found, originated as a qanat.
Japan
In Japan there are several dozen qanat-like structures, locally known as 'mambo' or 'manbo', most notably in the Mie and Gifu Prefectures. Whereas some link their origin clearly to the Chinese karez, and therefore to the Iranian source, a Japanese conference in 2008 found insufficient scientific studies to evaluate the origins of the mambo.
Jordan
Among the qanats built in the Roman Empire, the long Gadara Aqueduct in northern Jordan was possibly the longest continuous qanat ever built. Partly following the course of an older Hellenistic aqueduct, excavation work arguably started after a visit by emperor Hadrian in 129–130 CE. The Gadara Aqueduct was never totally finished and was put in service only in sections.
Pakistan
In Pakistan qanat irrigation system is endemic only in Balochistan. The major concentration is in the north and northwest along the Pakistan-Afghanistan border and oasis of Makoran division. The karez system of the Balochistan desert is on the tentative list for future world heritage sites in Pakistan.
The acute shortage of water resources give water a decisive role in the regional conflicts arose in the course of history of Balochistan. Therefore, in Balochistan, the possession of water resources is more important than ownership of land. Hence afterward a complex system for the collection, channeling and distribution of water was developed in Balochistan. Similarly, the distribution and unbiased flow of water to different stockholders also necessitate the importance of different societal classes in Balochistan in general and particularly in Makoran.
For instance, sarrishta, literally, head of the chain, is responsible for administration of channel. He normally owns the largest water quota. Under sarrishta, there are several heads of owners issadar who also possessed larger water quotas. The social hierarchy within Baloch society of Makoran depends upon the possession of largest quotas of water. The role of sarrishta in some cases hierarchical and passing from generations within the family and he must have the knowledge of the criteria of unbiased distribution of water among different issadar.
The sharing of water is based on a complex indigenous system of measurement depends upon time and space particularly to the phases of moon; the hangams. Based on seasonal variations and share of water the hangams are apportioned among various owners over period of seven or fourteen days. However, in some places, instead of hangam, anna used which is based on twelve-hour period for each quota. Therefore, if a person own 16 quotas it means that he is entitled for water for eight days in high seasons and 16 days in winter when water level went down as well as expectation of winter rain (Baharga) in Makran region. The twelve-hour water quota again subdivided into several sub-fractions of local measuring scales such as tas or pad (Dr Gul Hasan Pro VC LUAWMS, 2 day National conference on Kech).
The Chagai district is in the north west corner of Balochistan, Pakistan, bordering with Afghanistan and Iran. Qanats, locally known as Kahn, are found more broadly in this region. They are spread from Chaghai district all the way up to Zhob district.
Syria
Qanats were found over much of Syria. The widespread installation of groundwater pumps has lowered the water table and qanat system. Qanats have gone dry and been abandoned across the country.
Oman
In Oman from the Iron Age period (found in Salut, Bat and other sites) a system of underground aqueducts called 'Falaj' were constructed, a series of well-like vertical shafts, connected by gently sloping horizontal tunnels. There are three types of Falaj: Daudi () with underground aqueducts, Ghaili () requiring a dam to collect the water, and Aini () whose source is a water spring. These enabled large scale agriculture to flourish in a dryland environment. According to UNESCO, some 3,000 aflaj (plural) or falaj (singular), are still in use in Oman today. Nizwa, the former capital city of Oman, was built around a falaj which is in use to this day. These systems date to before the Iron Age in Oman. In July 2006, five representative examples of this irrigation system were inscribed as a World Heritage Site.
United Arab Emirates
The oases of the city of Al Ain (particularly Al-Ain, Al-Qattarah, Al-Mu'taredh, Al-Jimi, Al-Muwaiji, and Hili), adjacent to Al-Buraimi in Oman, continue traditional falaj (qanat) irrigations for the palm groves and gardens, and form part of the city's ancient heritage. Multiple aflaj have been found from the early Iron Age, as early as 1100 BC, and some sources claim these to be the earliest examples of the qanat irrigation system.
The falaj system continued to be in use into the early Pre Islamic (300 BC - 300 AD) and were reintroduced after the early Islamic conquests in the souqs of Julfar, Dibba, and Tawwam. Islamic geographer Al Muqqadasi stated: "Hafit {Tuwwam} abounds in palm trees" in the 10th century, indicating extensive use of the falaj system. The falaj system is still operating in the city of Al Ain, as well as in multiple mountainous settlements, including the villages of Wadi Shees and Masafi.
Europe
Greece
The Tunnel of Eupalinos on Samos runs for 1 kilometre through a hill to supply water to Pythagorion. It was built on the order of Polycrates around 550 BCE. At either end of the tunnel proper, shallow qanat-like tunnels carried the water from the spring and to the town.
Italy
The long Tunnels of Claudius, intended to partially drain the largest Italian inland water, Fucine Lake, was constructed using the qanat technique. It featured shafts up to 122 m deep. The entire ancient town of Palermo in Sicily was equipped with a huge qanat system built during the Arab period (827–1072). Many of the qanats are now mapped and some can be visited. The famous Scirocco room has an air-conditioning system cooled by the flow of water in a qanat and a "wind tower", a structure able to catch the wind and use it to draw the cooled air up into the room.
Luxembourg
The Raschpëtzer near Helmsange in southern Luxembourg is a particularly well preserved example of a Roman qanat. It is probably the most extensive system of its kind north of the Alps. To date, some 330 m of the total tunnel length of 600 m have been explored. Thirteen of the 20 to 25 shafts have been investigated. The qanat appears to have provided water for a large Roman villa on the slopes of the Alzette valley. It was built during the Gallo-Roman period, probably around the year 150 and functioned for about 120 years thereafter.
Spain
There are still many examples of or qanat systems in Spain, most likely brought to the area by the Moors during their rule of the Iberian peninsula. Turrillas in Andalusia on the north facing slopes of the Sierra de Alhamilla has evidence of a qanat system. Granada is another site with an extensive qanat system. In Madrid they were called and were used until the construction of the Canal de Isabel II. See and in Spanish.
The Americas
Qanats in the Americas, usually referred to as puquios or filtration galleries, can be found in the Nazca Province of Peru and in northern Chile. The origin and dating of the Nazca puquios is disputed, although some archaeologists have asserted that they were constructed by the indigenous people of the Nazca culture beginning about 500 CE.
The Spanish introduced qanats into Mexico in 1520 CE.
In the Atacama Desert of northern Chile the shafts of puquios are known as socavones. Socavones are known to exist in Azapa Valley and the oasis of Sibaya, Pica-Matilla, and Puquio de Núñez. In 1918 geologist Juan Brüggen mentioned the existence of 23 socavones in the Pica oasis, yet these have since then been abandoned due to economic and social changes.
Symbolism in Iranian culture
In an August 21, 1906, letter written from Tehran, Florence Khanum, the American wife of Persian diplomat Ali Kuli Khan, described the use of qanats for the garden at the home of her brother-in-law, General Husayn Kalantar,
January 1, 1913
An old tradition in Iran was to hold symbolic wedding ceremonies between widows and qanats in which the widow became the "wife" of the qanat. This was believed to help ensure the continued flow of water.
| Technology | Food, water and health | null |
536063 | https://en.wikipedia.org/wiki/Exothermic%20reaction | Exothermic reaction | In thermochemistry, an exothermic reaction is a "reaction for which the overall standard enthalpy change ΔH⚬ is negative." Exothermic reactions usually release heat. The term is often confused with exergonic reaction, which IUPAC defines as "... a reaction for which the overall standard Gibbs energy change ΔG⚬ is negative." A strongly exothermic reaction will usually also be exergonic because ΔH⚬ makes a major contribution to ΔG⚬. Most of the spectacular chemical reactions that are demonstrated in classrooms are exothermic and exergonic. The opposite is an endothermic reaction, which usually takes up heat and is driven by an entropy increase in the system.
Examples
Examples are numerous: combustion, the thermite reaction, combining strong acids and bases, polymerizations. As an example in everyday life, hand warmers make use of the oxidation of iron to achieve an exothermic reaction:
4Fe + 3O2 → 2Fe2O3 ΔH⚬ = - 1648 kJ/mol
A particularly important class of exothermic reactions is combustion of a hydrocarbon fuel, e.g. the burning of natural gas:
CH4 + 2O2 → CO2 + 2H2O ΔH⚬ = - 890 kJ/mol
These sample reactions are strongly exothermic.
Uncontrolled exothermic reactions, those leading to fires and explosions, are wasteful because it is difficult to capture the released energy. Nature effects combustion reactions under highly controlled conditions, avoiding fires and explosions, in aerobic respiration so as to capture the released energy, e.g. for the formation of ATP.
Measurement
The enthalpy of a chemical system is essentially its energy. The enthalpy change ΔH for a reaction is equal to the heat q transferred out of (or into) a closed system at constant pressure without in- or output of electrical energy. Heat production or absorption in a chemical reaction is measured using calorimetry, e.g. with a bomb calorimeter. One common laboratory instrument is the reaction calorimeter, where the heat flow from or into the reaction vessel is monitored. The heat release and corresponding energy change, Δ, of a combustion reaction can be measured particularly accurately.
The measured heat energy released in an exothermic reaction is converted to ΔH⚬ in Joule per mole (formerly cal/mol). The standard enthalpy change ΔH⚬ is essentially the enthalpy change when the stoichiometric coefficients in the reaction are considered as the amounts of reactants and products (in mole); usually, the initial and final temperature is assumed to be 25 °C. For gas-phase reactions, ΔH⚬ values are related to bond energies to a good approximation by:
Δ⚬ = total bond energy of reactants − total bond energy of products
In an exothermic reaction, by definition, the enthalpy change has a negative value:
Δ = Hproducts - Hreactants < 0
where a larger value (the higher energy of the reactants) is subtracted from a smaller value (the lower energy of the products). For example, when hydrogen burns:
2H2 (g) + O2 (g) → 2H2O (g)
Δ⚬ = −483.6 kJ/mol
| Physical sciences | Thermodynamics | Chemistry |
536079 | https://en.wikipedia.org/wiki/Unconformity | Unconformity | An unconformity is a buried erosional or non-depositional surface separating two rock masses or strata of different ages, indicating that sediment deposition was not continuous. In general, the older layer was exposed to erosion for an interval of time before deposition of the younger layer, but the term is used to describe any break in the sedimentary geologic record. The significance of angular unconformity (see below) was shown by James Hutton, who found examples of Hutton's Unconformity at Jedburgh in 1787 and at Siccar Point in Berwickshire in 1788, both in Scotland.
The rocks above an unconformity are younger than the rocks beneath (unless the sequence has been overturned). An unconformity represents time during which no sediments were preserved in a region or were subsequently eroded before the next deposition. The local record for that time interval is missing and geologists must use other clues to discover that part of the geologic history of that area. The interval of geologic time not represented is called a hiatus. It is a kind of relative dating.
Types
Disconformity
A disconformity is an unconformity between parallel layers of sedimentary rocks which represents a period of erosion or non-deposition. Disconformities are marked by features of subaerial erosion. This type of erosion can leave channels and paleosols in the rock record.
Nonconformity
A nonconformity exists between sedimentary rocks and metamorphic or igneous rocks when the sedimentary rock lies above and was deposited on the pre-existing and eroded metamorphic or igneous rock. Namely, if the rock below the break is igneous or has lost its bedding due to metamorphism, then the plane of juncture is a nonconformity.
Angular unconformity
An angular unconformity is an unconformity where horizontally parallel strata of sedimentary rock are deposited on tilted and eroded layers, producing an angular discordance with the overlying horizontal layers. The whole sequence may later be deformed and tilted by further orogenic activity. A typical case history is presented by the Briançonnais realm (Swiss and French Prealps) during the Jurassic.
Angular unconformities can occur in ash fall layers of pyroclastic rock deposited by volcanoes during explosive eruptions. In these cases, the hiatus in deposition represented by the unconformity may be geologically very short – hours, days or weeks.
Paraconformity
A paraconformity is a type of unconformity in which the sedimentary layers above and below the unconformity are parallel, but there is no obvious erosional break between them. A break in sedimentation is indicated, for example, by fossil evidence. It is also called nondepositional unconformity or pseudoconformity. Short paraconformities are called diastems.
Buttress unconformity
A buttress unconformity, also known as onlap unconformity, occurs when younger bedding is deposited against older strata thus influencing its bedding structure.
Blended unconformity
A blended unconformity is a type of disconformity or nonconformity with no distinct separation plane or contact, sometimes consisting of soils, paleosols, or beds of pebbles derived from the underlying rock.
Gallery
| Physical sciences | Stratigraphy | Earth science |
536080 | https://en.wikipedia.org/wiki/Student%27s%20t-test | Student's t-test | Student's t-test is a statistical test used to test whether the difference between the response of two groups is statistically significant or not. It is any statistical hypothesis test in which the test statistic follows a Student's t-distribution under the null hypothesis. It is most commonly applied when the test statistic would follow a normal distribution if the value of a scaling term in the test statistic were known (typically, the scaling term is unknown and is therefore a nuisance parameter). When the scaling term is estimated based on the data, the test statistic—under certain conditions—follows a Student's t distribution. The t-test's most common application is to test whether the means of two populations are significantly different. In many cases, a Z-test will yield very similar results to a t-test because the latter converges to the former as the size of the dataset increases.
History
The term "t-statistic" is abbreviated from "hypothesis test statistic". In statistics, the t-distribution was first derived as a posterior distribution in 1876 by Helmert and Lüroth. The t-distribution also appeared in a more general form as Pearson type IV distribution in Karl Pearson's 1895 paper. However, the t-distribution, also known as Student's t-distribution, gets its name from William Sealy Gosset, who first published it in English in 1908 in the scientific journal Biometrika using the pseudonym "Student" because his employer preferred staff to use pen names when publishing scientific papers. Gosset worked at the Guinness Brewery in Dublin, Ireland, and was interested in the problems of small samples for example, the chemical properties of barley with small sample sizes. Hence a second version of the etymology of the term Student is that Guinness did not want their competitors to know that they were using the t-test to determine the quality of raw material. Although it was William Gosset after whom the term "Student" is penned, it was actually through the work of Ronald Fisher that the distribution became well known as "Student's distribution" and "Student's t-test".
Gosset devised the t-test as an economical way to monitor the quality of stout. The t-test work was submitted to and accepted in the journal Biometrika and published in 1908.
Guinness had a policy of allowing technical staff leave for study (so-called "study leave"), which Gosset used during the first two terms of the 1906–1907 academic year in Professor Karl Pearson's Biometric Laboratory at University College London. Gosset's identity was then known to fellow statisticians and to editor-in-chief Karl Pearson.
Uses
One-sample t-test
A one-sample Student's t-test is a location test of whether the mean of a population has a value specified in a null hypothesis. In testing the null hypothesis that the population mean is equal to a specified value , one uses the statistic
where is the sample mean, is the sample standard deviation and is the sample size. The degrees of freedom used in this test are . Although the parent population does not need to be normally distributed, the distribution of the population of sample means is assumed to be normal.
By the central limit theorem, if the observations are independent and the second moment exists, then will be approximately normal .
Two-sample t-tests
A two-sample location test of the null hypothesis such that the means of two populations are equal. All such tests are usually called Student's t-tests, though strictly speaking that name should only be used if the variances of the two populations are also assumed to be equal; the form of the test used when this assumption is dropped is sometimes called Welch's t-test. These tests are often referred to as unpaired or independent samples t-tests, as they are typically applied when the statistical units underlying the two samples being compared are non-overlapping.
Two-sample t-tests for a difference in means involve independent samples (unpaired samples) or paired samples. Paired t-tests are a form of blocking, and have greater power (probability of avoiding a type II error, also known as a false negative) than unpaired tests when the paired units are similar with respect to "noise factors" (see confounder) that are independent of membership in the two groups being compared. In a different context, paired t-tests can be used to reduce the effects of confounding factors in an observational study.
Independent (unpaired) samples
The independent samples t-test is used when two separate sets of independent and identically distributed samples are obtained, and one variable from each of the two populations is compared. For example, suppose we are evaluating the effect of a medical treatment, and we enroll 100 subjects into our study, then randomly assign 50 subjects to the treatment group and 50 subjects to the control group. In this case, we have two independent samples and would use the unpaired form of the t-test.
Paired samples
Paired samples t-tests typically consist of a sample of matched pairs of similar units, or one group of units that has been tested twice (a "repeated measures" t-test).
A typical example of the repeated measures t-test would be where subjects are tested prior to a treatment, say for high blood pressure, and the same subjects are tested again after treatment with a blood-pressure-lowering medication. By comparing the same patient's numbers before and after treatment, we are effectively using each patient as their own control. That way the correct rejection of the null hypothesis (here: of no difference made by the treatment) can become much more likely, with statistical power increasing simply because the random interpatient variation has now been eliminated. However, an increase of statistical power comes at a price: more tests are required, each subject having to be tested twice. Because half of the sample now depends on the other half, the paired version of Student's t-test has only degrees of freedom (with being the total number of observations). Pairs become individual test units, and the sample has to be doubled to achieve the same number of degrees of freedom. Normally, there are degrees of freedom (with being the total number of observations).
A paired samples t-test based on a "matched-pairs sample" results from an unpaired sample that is subsequently used to form a paired sample, by using additional variables that were measured along with the variable of interest. The matching is carried out by identifying pairs of values consisting of one observation from each of the two samples, where the pair is similar in terms of other measured variables. This approach is sometimes used in observational studies to reduce or eliminate the effects of confounding factors.
Paired samples t-tests are often referred to as "dependent samples t-tests".
Assumptions
Most test statistics have the form , where and are functions of the data.
may be sensitive to the alternative hypothesis (i.e., its magnitude tends to be larger when the alternative hypothesis is true), whereas is a scaling parameter that allows the distribution of to be determined.
As an example, in the one-sample t-test
where is the sample mean from a sample , of size , is the standard error of the mean, is the estimate of the standard deviation of the population, and is the population mean.
The assumptions underlying a t-test in the simplest form above are that:
follows a normal distribution with mean and variance .
follows a distribution with degrees of freedom. This assumption is met when the observations used for estimating come from a normal distribution (and i.i.d. for each group).
and are independent.
In the t-test comparing the means of two independent samples, the following assumptions should be met:
The means of the two populations being compared should follow normal distributions. Under weak assumptions, this follows in large samples from the central limit theorem, even when the distribution of observations in each group is non-normal.
If using Student's original definition of the t-test, the two populations being compared should have the same variance (testable using F-test, Levene's test, Bartlett's test, or the Brown–Forsythe test; or assessable graphically using a Q–Q plot). If the sample sizes in the two groups being compared are equal, Student's original t-test is highly robust to the presence of unequal variances. Welch's t-test is insensitive to equality of the variances regardless of whether the sample sizes are similar.
The data used to carry out the test should either be sampled independently from the two populations being compared or be fully paired. This is in general not testable from the data, but if the data are known to be dependent (e.g. paired by test design), a dependent test has to be applied. For partially paired data, the classical independent t-tests may give invalid results as the test statistic might not follow a t distribution, while the dependent t-test is sub-optimal as it discards the unpaired data.
Most two-sample t-tests are robust to all but large deviations from the assumptions.
For exactness, the t-test and Z-test require normality of the sample means, and the t-test additionally requires that the sample variance follows a scaled χ distribution, and that the sample mean and sample variance be statistically independent. Normality of the individual data values is not required if these conditions are met. By the central limit theorem, sample means of moderately large samples are often well-approximated by a normal distribution even if the data are not normally distributed. However, the sample size required for the sample means to converge to normality depends on the skewness of the distribution of the original data. The sample can vary from 30 to 100 or higher values depending on the skewness. F
For non-normal data, the distribution of the sample variance may deviate substantially from a χ distribution.
However, if the sample size is large, Slutsky's theorem implies that the distribution of the sample variance has little effect on the distribution of the test statistic. That is, as sample size increases:
as per the Central limit theorem,
as per the law of large numbers,
.
Calculations
Explicit expressions that can be used to carry out various t-tests are given below. In each case, the formula for a test statistic that either exactly follows or closely approximates a t-distribution under the null hypothesis is given. Also, the appropriate degrees of freedom are given in each case. Each of these statistics can be used to carry out either a one-tailed or two-tailed test.
Once the t value and degrees of freedom are determined, a p-value can be found using a table of values from Student's t-distribution. If the calculated p-value is below the threshold chosen for statistical significance (usually the 0.10, the 0.05, or 0.01 level), then the null hypothesis is rejected in favor of the alternative hypothesis.
Slope of a regression line
Suppose one is fitting the model
where is known, and are unknown, is a normally distributed random variable with mean 0 and unknown variance , and is the outcome of interest. We want to test the null hypothesis that the slope is equal to some specified value (often taken to be 0, in which case the null hypothesis is that and are uncorrelated).
Let
Then
has a t-distribution with degrees of freedom if the null hypothesis is true. The standard error of the slope coefficient:
can be written in terms of the residuals. Let
Then score is given by
Another way to determine the score is
where r is the Pearson correlation coefficient.
The score, intercept can be determined from the score, slope:
where is the sample variance.
Independent two-sample t-test
Equal sample sizes and variance
Given two groups (1, 2), this test is only applicable when:
the two sample sizes are equal,
it can be assumed that the two distributions have the same variance.
Violations of these assumptions are discussed below.
The statistic to test whether the means are different can be calculated as follows:
where
Here is the pooled standard deviation for , and and are the unbiased estimators of the population variance. The denominator of is the standard error of the difference between two means.
For significance testing, the degrees of freedom for this test is , where is sample size.
Equal or unequal sample sizes, similar variances ( < < 2)
This test is used only when it can be assumed that the two distributions have the same variance (when this assumption is violated, see below).
The previous formulae are a special case of the formulae below, one recovers them when both samples are equal in size: .
The statistic to test whether the means are different can be calculated as follows:
where
is the pooled standard deviation of the two samples: it is defined in this way so that its square is an unbiased estimator of the common variance, whether or not the population means are the same. In these formulae, is the number of degrees of freedom for each group, and the total sample size minus two (that is, ) is the total number of degrees of freedom, which is used in significance testing.
The minimum detectable effect (MDE) is:
Equal or unequal sample sizes, unequal variances (sX1 > 2sX2 or sX2 > 2sX1)
This test, also known as Welch's t-test, is used only when the two population variances are not assumed to be equal (the two sample sizes may or may not be equal) and hence must be estimated separately. The statistic to test whether the population means are different is calculated as
where
Here is the unbiased estimator of the variance of each of the two samples with = number of participants in group ( = 1 or 2). In this case
is not a pooled variance. For use in significance testing, the distribution of the test statistic is approximated as an ordinary Student's t-distribution with the degrees of freedom calculated using
This is known as the Welch–Satterthwaite equation. The true distribution of the test statistic actually depends (slightly) on the two unknown population variances (see Behrens–Fisher problem).
Exact method for unequal variances and sample sizes
The test deals with the famous Behrens–Fisher problem, i.e., comparing the difference between the means of two normally distributed populations when the variances of the two populations are not assumed to be equal, based on two independent samples.
The test is developed as an exact test that allows for unequal sample sizes and unequal variances of two populations. The exact property still holds even with small extremely small and unbalanced sample sizes (e.g. ).
The statistic to test whether the means are different can be calculated as follows:
Let and be the i.i.d. sample vectors () from and separately.
Let be an orthogonal matrix whose elements of the first row are all , similarly, let be the first n rows of an orthogonal matrix (whose elements of the first row are all ).
Then is an n-dimensional normal random vector.
From the above distribution we see that
Dependent t-test for paired samples
This test is used when the samples are dependent; that is, when there is only one sample that has been tested twice (repeated measures) or when there are two samples that have been matched or "paired". This is an example of a paired difference test. The t statistic is calculated as
where and are the average and standard deviation of the differences between all pairs. The pairs are e.g. either one person's pre-test and post-test scores or between-pairs of persons matched into meaningful groups (for instance, drawn from the same family or age group: see table). The constant is zero if we want to test whether the average of the difference is significantly different. The degree of freedom used is , where represents the number of pairs.
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Worked examples
Let denote a set obtained by drawing a random sample of six measurements:
and let denote a second set obtained similarly:
These could be, for example, the weights of screws that were manufactured by two different machines.
We will carry out tests of the null hypothesis that the means of the populations from which the two samples were taken are equal.
The difference between the two sample means, each denoted by , which appears in the numerator for all the two-sample testing approaches discussed above, is
The sample standard deviations for the two samples are approximately 0.05 and 0.11, respectively. For such small samples, a test of equality between the two population variances would not be very powerful. Since the sample sizes are equal, the two forms of the two-sample t-test will perform similarly in this example.
Unequal variances
If the approach for unequal variances (discussed above) is followed, the results are
and the degrees of freedom
The test statistic is approximately 1.959, which gives a two-tailed test p-value of 0.09077.
Equal variances
If the approach for equal variances (discussed above) is followed, the results are
and the degrees of freedom
The test statistic is approximately equal to 1.959, which gives a two-tailed p-value of 0.07857.
Related statistical tests
Alternatives to the t-test for location problems
The t-test provides an exact test for the equality of the means of two i.i.d. normal populations with unknown, but equal, variances. (Welch's t-test is a nearly exact test for the case where the data are normal but the variances may differ.) For moderately large samples and a one tailed test, the t-test is relatively robust to moderate violations of the normality assumption. In large enough samples, the t-test asymptotically approaches the z-test, and becomes robust even to large deviations from normality.
If the data are substantially non-normal and the sample size is small, the t-test can give misleading results. See Location test for Gaussian scale mixture distributions for some theory related to one particular family of non-normal distributions.
When the normality assumption does not hold, a non-parametric alternative to the t-test may have better statistical power. However, when data are non-normal with differing variances between groups, a t-test may have better type-1 error control than some non-parametric alternatives. Furthermore, non-parametric methods, such as the Mann-Whitney U test discussed below, typically do not test for a difference of means, so should be used carefully if a difference of means is of primary scientific interest. For example, Mann-Whitney U test will keep the type 1 error at the desired level alpha if both groups have the same distribution. It will also have power in detecting an alternative by which group B has the same distribution as A but after some shift by a constant (in which case there would indeed be a difference in the means of the two groups). However, there could be cases where group A and B will have different distributions but with the same means (such as two distributions, one with positive skewness and the other with a negative one, but shifted so to have the same means). In such cases, MW could have more than alpha level power in rejecting the Null hypothesis but attributing the interpretation of difference in means to such a result would be incorrect.
In the presence of an outlier, the t-test is not robust. For example, for two independent samples when the data distributions are asymmetric (that is, the distributions are skewed) or the distributions have large tails, then the Wilcoxon rank-sum test (also known as the Mann–Whitney U test) can have three to four times higher power than the t-test. The nonparametric counterpart to the paired samples t-test is the Wilcoxon signed-rank test for paired samples. For a discussion on choosing between the t-test and nonparametric alternatives, see Lumley, et al. (2002).
One-way analysis of variance (ANOVA) generalizes the two-sample t-test when the data belong to more than two groups.
A design which includes both paired observations and independent observations
When both paired observations and independent observations are present in the two sample design, assuming data are missing completely at random (MCAR), the paired observations or independent observations may be discarded in order to proceed with the standard tests above. Alternatively making use of all of the available data, assuming normality and MCAR, the generalized partially overlapping samples t-test could be used.
Multivariate testing
A generalization of Student's t statistic, called Hotelling's t-squared statistic, allows for the testing of hypotheses on multiple (often correlated) measures within the same sample. For instance, a researcher might submit a number of subjects to a personality test consisting of multiple personality scales (e.g. the Minnesota Multiphasic Personality Inventory). Because measures of this type are usually positively correlated, it is not advisable to conduct separate univariate t-tests to test hypotheses, as these would neglect the covariance among measures and inflate the chance of falsely rejecting at least one hypothesis (Type I error). In this case a single multivariate test is preferable for hypothesis testing. Fisher's Method for combining multiple tests with alpha reduced for positive correlation among tests is one. Another is Hotelling's T statistic follows a T distribution. However, in practice the distribution is rarely used, since tabulated values for T are hard to find. Usually, T is converted instead to an F statistic.
For a one-sample multivariate test, the hypothesis is that the mean vector () is equal to a given vector (). The test statistic is Hotelling's t:
where is the sample size, is the vector of column means and is an sample covariance matrix.
For a two-sample multivariate test, the hypothesis is that the mean vectors () of two samples are equal. The test statistic is Hotelling's two-sample t:
The two-sample t-test is a special case of simple linear regression
The two-sample t-test is a special case of simple linear regression as illustrated by the following example.
A clinical trial examines 6 patients given drug or placebo. Three (3) patients get 0 units of drug (the placebo group). Three (3) patients get 1 unit of drug (the active treatment group). At the end of treatment, the researchers measure the change from baseline in the number of words that each patient can recall in a memory test.
A table of the patients' word recall and drug dose values are shown below.
Data and code are given for the analysis using the R programming language with the t.test and lmfunctions for the t-test and linear regression. Here are the same (fictitious) data above generated in R.
> word.recall.data=data.frame(drug.dose=c(0,0,0,1,1,1), word.recall=c(1,2,3,5,6,7))
Perform the t-test. Notice that the assumption of equal variance, var.equal=T, is required to make the analysis exactly equivalent to simple linear regression.
> with(word.recall.data, t.test(word.recall~drug.dose, var.equal=T))
Running the R code gives the following results.
The mean word.recall in the 0 drug.dose group is 2.
The mean word.recall in the 1 drug.dose group is 6.
The difference between treatment groups in the mean word.recall is 6 – 2 = 4.
The difference in word.recall between drug doses is significant (p=0.00805).
Perform a linear regression of the same data. Calculations may be performed using the R function lm() for a linear model.
> word.recall.data.lm = lm(word.recall~drug.dose, data=word.recall.data)
> summary(word.recall.data.lm)
The linear regression provides a table of coefficients and p-values.
The table of coefficients gives the following results.
The estimate value of 2 for the intercept is the mean value of the word recall when the drug dose is 0.
The estimate value of 4 for the drug dose indicates that for a 1-unit change in drug dose (from 0 to 1) there is a 4-unit change in mean word recall (from 2 to 6). This is the slope of the line joining the two group means.
The p-value that the slope of 4 is different from 0 is p = 0.00805.
The coefficients for the linear regression specify the slope and intercept of the line that joins the two group means, as illustrated in the graph. The intercept is 2 and the slope is 4.
Compare the result from the linear regression to the result from the t-test.
From the t-test, the difference between the group means is 6-2=4.
From the regression, the slope is also 4 indicating that a 1-unit change in drug dose (from 0 to 1) gives a 4-unit change in mean word recall (from 2 to 6).
The t-test p-value for the difference in means, and the regression p-value for the slope, are both 0.00805. The methods give identical results.
This example shows that, for the special case of a simple linear regression where there is a single x-variable that has values 0 and 1, the t-test gives the same results as the linear regression. The relationship can also be shown algebraically.
Recognizing this relationship between the t-test and linear regression facilitates the use of multiple linear regression and multi-way analysis of variance. These alternatives to t-tests allow for the inclusion of additional explanatory variables that are associated with the response. Including such additional explanatory variables using regression or anova reduces the otherwise unexplained variance, and commonly yields greater power to detect differences than do two-sample t-tests.
Software implementations
Many spreadsheet programs and statistics packages, such as QtiPlot, LibreOffice Calc, Microsoft Excel, SAS, SPSS, Stata, DAP, gretl, R, Python, PSPP, Wolfram Mathematica, MATLAB and Minitab, include implementations of Student's t-test.
| Mathematics | Statistics and probability | null |
536158 | https://en.wikipedia.org/wiki/Opisthokont | Opisthokont | The opisthokonts () are a broad group of eukaryotes, including both the animal and fungus kingdoms. The opisthokonts, previously called the "Fungi/Metazoa group", are generally recognized as a clade. Opisthokonts together with Apusomonadida and Breviata comprise the larger clade Obazoa.
Flagella and other characteristics
A common characteristic of opisthokonts is that flagellate cells, such as the sperm of most animals and the spores of the chytrid fungi, propel themselves with a single posterior flagellum. It is this feature that gives the group its name. In contrast, flagellate cells in other eukaryote groups propel themselves with one or more anterior flagella. Flagellate cells however have been secondarily lost in some opisthokont groups, including most of the fungi.
Opisthokont characteristics include synthesis of extracellular chitin in exoskeleton, cyst/spore wall, or cell wall of filamentous growth and hyphae; the extracellular digestion of substrates with osmotrophic absorption of nutrients; and other cell biosynthetic and metabolic pathways. Genera at the base of each clade are amoeboid and phagotrophic.
History
The close relationship between animals and fungi was suggested by Thomas Cavalier-Smith in 1987, who used the informal name opisthokonta (the formal name has been used for the chytrids by Copeland in 1956), and was supported by later genetic studies.
Early phylogenies placed fungi near the plants and other groups that have mitochondria with flat cristae, but this character varies. More recently, it has been said that holozoa (animals) and holomycota (fungi) are much more closely related to each other than either is to plants, because opisthokonts have a triple fusion of carbamoyl phosphate synthetase, dihydroorotase, and aspartate carbamoyltransferase that is not present in plants, and plants have a fusion of thymidylate synthase and dihydrofolate reductase not present in the opisthokonts. Animals and fungi are also more closely related to amoebas than to plants, and plants are more closely related to the SAR supergroup of protists than to animals or fungi. Animals and fungi are both heterotrophs, unlike plants, and while fungi are sessile like plants, there are also sessile animals.
Cavalier-Smith and Stechmann argue that the uniciliate eukaryotes such as opisthokonts and Amoebozoa, collectively called unikonts, split off from the other biciliate eukaryotes, called bikonts, shortly after they evolved.
Taxonomy
Opisthokonts are divided into Holomycota or Nucletmycea (fungi and all organisms more closely related to fungi than to animals) and Holozoa (animals and all organisms more closely related to animals than to fungi); no opisthokonts basal to the Holomycota/Holozoa split have yet been identified. The Opisthokonts was largely resolved by Torriella et al. Holomycota and Holozoa are composed of the following groups.
Holomycota (Fungus-like)
Fungi
Includes:
chytrids (flagellated, zoosporic fungi)
Fonticula (more recent work considers this to be part of Cristidiscoidea, a sister group to the fungi)
Hyaloraphidium (previously thought to be a green alga, now considered a fungus)
microsporidia (previously thought to be apicomplexia)
Nucleariida (more recent work considers this to be part of Cristidiscoidea, a sister group to the fungi)
Excludes:
labyrinthulomycetes (slime nets) (now included in the SAR supergroup)
myxomycetes (now included in amoebozoans)
oomycetes (water molds) (now included in the SAR supergroup)
Rozellida (placement uncertain)
Holozoa (Animal-like)
Corallochytrium (formerly considered a Heterokont)
Filozoa
Animalia (including myxozoa)
Choanoflagellata (flagellates formerly included in protozoa)
Filasterea
Mesomycetozoea
Amoebidiales (formerly considered trichomycetes)
Dermocystida (formerly considered parasitic fungi or sporozoans)
Eccrinales (formerly considered fungi)
Ichthyophonida (formerly considered parasitic fungi incertae sedis)
Phylogeny
The following phylogenetic tree indicates the evolutionary relationships between the different opisthokont lineages, and the time divergence of the clades in millions of years ago (Mya).
Gallery
| Biology and health sciences | Eukaryotes | Plants |
536313 | https://en.wikipedia.org/wiki/Polycarbonate | Polycarbonate | Polycarbonates (PC) are a group of thermoplastic polymers containing carbonate groups in their chemical structures. Polycarbonates used in engineering are strong, tough materials, and some grades are optically transparent. They are easily worked, molded, and thermoformed. Because of these properties, polycarbonates find many applications. Polycarbonates do not have a unique resin identification code (RIC) and are identified as "Other", 7 on the RIC list. Products made from polycarbonate can contain the precursor monomer bisphenol A (BPA).
Structure
Carbonate esters have planar OC(OC)2 cores, which confer rigidity. The unique O=C bond is short (1.173 Å in the depicted example), while the C-O bonds are more ether-like (the bond distances of 1.326 Å for the example depicted). Polycarbonates received their name because they are polymers containing carbonate groups (−O−(C=O)−O−). A balance of useful features, including temperature resistance, impact resistance and optical properties, positions polycarbonates between commodity plastics and engineering plastics.
Production
Phosgene route
The main polycarbonate material is produced by the reaction of bisphenol A (BPA) and phosgene . The overall reaction can be written as follows:
The first step of the synthesis involves treatment of bisphenol A with sodium hydroxide, which deprotonates the hydroxyl groups of the bisphenol A.
(HOC6H4)2CMe2 + 2 NaOH → Na2(OC6H4)2CMe2 + 2 H2O
The diphenoxide (Na2(OC6H4)2CMe2) reacts with phosgene to give a chloroformate, which subsequently is attacked by another phenoxide. The net reaction from the diphenoxide is:
Na2(OC6H4)2CMe2 + COCl2 → 1/n [OC(OC6H4)2CMe2]n + 2 NaCl
In this way, approximately one billion kilograms of polycarbonate is produced annually. Many other diols have been tested in place of bisphenol A, e.g. 1,1-bis(4-hydroxyphenyl)cyclohexane and dihydroxybenzophenone. The cyclohexane is used as a comonomer to suppress crystallisation tendency of the BPA-derived product. Tetrabromobisphenol A is used to enhance fire resistance. Tetramethylcyclobutanediol has been developed as a replacement for BPA.
Transesterification route
An alternative route to polycarbonates entails transesterification from BPA and diphenyl carbonate:
(HOC6H4)2CMe2 + (C6H5O)2CO → 1/n [OC(OC6H4)2CMe2]n + 2 C6H5OH
Properties and processing
Polycarbonate is a durable material. Although it has high impact-resistance, it has low scratch-resistance. Therefore, a hard coating is applied to polycarbonate eyewear lenses and polycarbonate exterior automotive components. The characteristics of polycarbonate compare to those of polymethyl methacrylate (PMMA, acrylic), but polycarbonate is stronger and will hold up longer to extreme temperature. Thermally processed material is usually totally amorphous, and as a result is highly transparent to visible light, with better light transmission than many kinds of glass.
Polycarbonate has a glass transition temperature of about , so it softens gradually above this point and flows above about . Tools must be held at high temperatures, generally above to make strain-free and stress-free products. Low molecular mass grades are easier to mold than higher grades, but their strength is lower as a result. The toughest grades have the highest molecular mass, but are more difficult to process.
Unlike most thermoplastics, polycarbonate can undergo large plastic deformations without cracking or breaking. As a result, it can be processed and formed at room temperature using sheet metal techniques, such as bending on a brake. Even for sharp angle bends with a tight radius, heating may not be necessary. This makes it valuable in prototyping applications where transparent or electrically non-conductive parts are needed, which cannot be made from sheet metal. PMMA/Acrylic, which is similar in appearance to polycarbonate, is brittle and cannot be bent at room temperature.
Main transformation techniques for polycarbonate resins:
extrusion into tubes, rods and other profiles including multiwall
extrusion with cylinders (calenders) into sheets () and films (below ), which can be used directly or manufactured into other shapes using thermoforming or secondary fabrication techniques, such as bending, drilling, or routing. Due to its chemical properties it is not conducive to laser-cutting.
injection molding into ready articles
Polycarbonate may become brittle when exposed to ionizing radiation above
Applications
Electronic components
Polycarbonate is mainly used for electronic applications that capitalize on its collective safety features. A good electrical insulator with heat-resistant and flame-retardant properties, it is used in products associated with power systems and telecommunications hardware. It can serve as a dielectric in high-stability capacitors. Commercial manufacture of polycarbonate capacitors mostly stopped after sole manufacturer Bayer AG stopped making capacitor-grade polycarbonate film at the end of 2000.
Construction materials
The second largest consumer of polycarbonates is the construction industry, e.g. for domelights, flat or curved glazing, roofing sheets and sound walls.
Polycarbonates are used to create materials used in buildings that must be durable but light.
3D printing
Polycarbonates are used extensively in 3D FDM printing, producing durable strong plastic products with a high melting point. Polycarbonate is relatively difficult for casual hobbyists to print compared to thermoplastics such as Polylactic acid (PLA) or Acrylonitrile butadiene styrene (ABS) because of the high melting point, difficulty with print bed adhesion, tendency to warp during printing, and tendency to absorb moisture in humid environments. Despite these issues, 3D printing using polycarbonates is common in the professional community.
Data storage
A major polycarbonate market is the production of compact discs, DVDs, and Blu-ray discs. These discs are produced by injection-molding polycarbonate into a mold cavity that has on one side a metal stamper containing a negative image of the disc data, while the other mold side is a mirrored surface. Typical products of sheet/film production include applications in advertisement (signs, displays, poster protection).
Automotive, aircraft, and security components
In the automotive industry, injection-molded polycarbonate can produce very smooth surfaces that make it well-suited for sputter deposition or evaporation deposition of aluminium without the need for a base-coat. Decorative bezels and optical reflectors are commonly made of polycarbonate.
Its low weight and high impact resistance have made polycarbonate the dominant material for automotive headlamp lenses. However, automotive headlamps require outer surface coatings because of its low scratch resistance and susceptibility to ultraviolet degradation (yellowing). The use of polycarbonate in automotive applications is limited to low stress applications. Stress from fasteners, plastic welding and molding render polycarbonate susceptible to stress corrosion cracking when it comes in contact with certain accelerants such as salt water and plastisol. It can be laminated to make bullet-proof "glass", although "bullet-resistant" is more accurate for the thinner windows, such as are used in bullet-resistant windows in automobiles. The thicker barriers of transparent plastic used in teller's windows and barriers in banks are also polycarbonate.
So-called "theft-proof" large plastic packaging for smaller items, which cannot be opened by hand, is typically made from polycarbonate.
The cockpit canopy of the Lockheed Martin F-22 Raptor jet fighter is fabricated from high optical quality polycarbonate. It is the largest item of its type.
Niche applications
Polycarbonate, being a versatile material with attractive processing and physical properties, has attracted myriad smaller applications. The use of injection molded drinking bottles, glasses and food containers is common, but the use of BPA in the manufacture of polycarbonate has stirred concerns (see Potential hazards in food contact applications), leading to development and use of "BPA-free" plastics in various formulations.
Polycarbonate is commonly used in eye protection, as well as in other projectile-resistant viewing and lighting applications that would normally indicate the use of glass, but require much higher impact-resistance. Polycarbonate lenses also protect the eye from UV light. Many kinds of lenses are manufactured from polycarbonate, including automotive headlamp lenses, lighting lenses, sunglass/eyeglass lenses, camera lenses, swimming goggles and SCUBA masks, and safety glasses/goggles/visors including visors in sporting helmets/masks and police riot gear (helmet visors, riot shields, etc.). Windscreens in small motorized vehicles are commonly made of polycarbonate, such as for motorcycles, ATVs, golf carts, and small airplanes and helicopters.
The light weight of polycarbonate as opposed to glass has led to development of electronic display screens that replace glass with polycarbonate, for use in mobile and portable devices. Such displays include newer e-ink and some LCD screens, though CRT, plasma screen and other LCD technologies generally still require glass for its higher melting temperature and its ability to be etched in finer detail.
As more and more governments are restricting the use of glass in pubs and clubs due to the increased incidence of glassings, polycarbonate glasses are becoming popular for serving alcohol because of their strength, durability, and glass-like feel.
Other miscellaneous items include durable, lightweight luggage, MP3/digital audio player cases, ocarinas, computer cases, riot shields, instrument panels, tealight candle containers and food blender jars. Many toys and hobby items are made from polycarbonate parts, like fins, gyro mounts, and flybar locks in radio-controlled helicopters, and transparent LEGO (ABS is used for opaque pieces).
Standard polycarbonate resins are not suitable for long term exposure to UV radiation. To overcome this, the primary resin can have UV stabilisers added. These grades are sold as UV stabilized polycarbonate to injection moulding and extrusion companies. Other applications, including polycarbonate sheets, may have the anti-UV layer added as a special coating or a coextrusion for enhanced weathering resistance.
Polycarbonate is also used as a printing substrate for nameplate and other forms of industrial grade under printed products. The polycarbonate provides a barrier to wear, the elements, and fading.
Medical applications
Many polycarbonate grades are used in medical applications and comply with both ISO 10993-1 and USP Class VI standards (occasionally referred to as PC-ISO). Class VI is the most stringent of the six USP ratings. These grades can be sterilized using steam at 120 °C, gamma radiation, or by the ethylene oxide (EtO) method. Trinseo strictly limits all its plastics with regard to medical applications. Aliphatic polycarbonates have been developed with improved biocompatibility and degradability for nanomedicine applications.
Mobile phones
Some smartphone manufacturers use polycarbonate. Nokia used polycarbonate in their phones starting with the N9's unibody case in 2011. This practice continued with various phones in the Lumia series. Samsung started using polycarbonate with Galaxy S III's hyperglaze-branded removable battery cover in 2012. This practice continues with various phones in the Galaxy series. Apple started using polycarbonate with the iPhone 5C's unibody case in 2013.
Benefits over glass and metal back covers include durability against shattering (advantage over glass), bending and scratching (advantage over metal), shock absorption, low manufacturing costs, and no interference with radio signals and wireless charging (advantage over metal).
Polycarbonate back covers are available in glossy or matte surface textures.
History
Polycarbonates were first discovered in 1898 by Alfred Einhorn, a German scientist working at the University of Munich. However, after 30 years' laboratory research, this class of materials was abandoned without commercialization. Research resumed in 1953, when Hermann Schnell at Bayer in Uerdingen, Germany patented the first linear polycarbonate. The brand name "Makrolon" was registered in 1955.
Also in 1953, and one week after the invention at Bayer, Daniel Fox at General Electric (GE) in Pittsfield, Massachusetts, independently synthesized a branched polycarbonate. Both companies filed for U.S. patents in 1955, and agreed that the company lacking priority would be granted a license to the technology.
Patent priority was resolved in Bayer's favor, and Bayer began commercial production under the trade name Makrolon in 1958. GE began production under the name Lexan in 1960, creating the GE Plastics division in 1973.
After 1970, the original brownish polycarbonate tint was improved to "glass-clear".
Potential hazards in food contact applications
The use of polycarbonate containers for the purpose of food storage is controversial. The basis of this controversy is their hydrolysis (degradation by water, often referred to as leaching) occurring at high temperature, releases bisphenol A:
1/n [OC(OC6H4)2CMe2]n + H2O → (HOC6H4)2CMe2 + CO2
More than 100 studies have explored the bioactivity of bisphenol A derived from polycarbonates. Bisphenol A appeared to be released from polycarbonate animal cages into water at room temperature and it may have been responsible for enlargement of the reproductive organs of female mice. However, the animal cages used in the research were fabricated from industrial grade polycarbonate, rather than FDA food grade polycarbonate.
An analysis of the literature on bisphenol A leachate low-dose effects by vom Saal and Hughes published in August 2005 seems to have found a suggestive correlation between the source of funding and the conclusion drawn. Industry-funded studies tend to find no significant effects whereas government-funded studies tend to find significant effects.
Sodium hypochlorite bleach and other alkali cleaners catalyze the release of the bisphenol A from polycarbonate containers. Polycarbonate is incompatible with ammonia and acetone. Alcohol is a recommended organic solvent for cleaning grease and oils from polycarbonate.
Environmental impact
Disposal
Studies have shown that at temperatures above 70 °C, and high humidity, polycarbonate will hydrolyze to bisphenol A (BPA). After about 30 days at 85 °C/96% RH, surface crystals are formed which for 70% consisted of BPA. BPA is a compound that is currently on the list of potential environmental hazardous chemicals. It is on the watch list of many countries, such as United States and Germany.
-(-OC6H4)2C(CH3)2CO-)-n + H2O → (CH3)2C(C6H4OH)2 + CO2
The leaching of BPA from polycarbonate can also occur at environmental temperature and normal pH (in landfills).The amount of leaching increases as the polycarbonate parts get older. A study found that the decomposition of BPA in landfills (under anaerobic conditions) will not occur. It will therefore be persistent in landfills. Eventually, it will find its way into water bodies and contribute to aquatic pollution.
Photo-oxidation of polycarbonate
In the presence of UV light, oxidation of this polymer yields compounds such as ketones, phenols, o-phenoxybenzoic acid, benzyl alcohol and other unsaturated compounds. This has been suggested through kinetic and spectral studies. The yellow color formed after long exposure to sun can also be related to further oxidation of phenolic end group
(OC6H4)2C(CH3)2CO )n + O2 , R* → (OC6H4)2C(CH3CH2)CO)n
This product can be further oxidized to form smaller unsaturated compounds. This can proceed via two different pathways, the products formed depends on which mechanism takes place.
Pathway A
(OC6H4)2C(CH3CH2)CO + O2, H* HO(OC6H4)OCO + CH3COCH2(OC6H4)OCO
Pathway B
(OC6H4)2C(CH3CH2)CO)n + O2, H* OCO(OC6H4)CH2OH + OCO(OC6H4)COCH3
Photo-aging reaction
Photo-aging is another degradation route for polycarbonates. Polycarbonate molecules (such as the aromatic ring) absorb UV radiation. This absorbed energy causes cleavage of covalent bonds which initiates the photo-aging process. The reaction can be propagated via side chain oxidation, ring oxidation or photo-Fries rearrangement. Products formed include phenyl salicylate, dihydroxybenzophenone groups, and hydroxydiphenyl ether groups.
(C16H14O3)n C16H17O3 + C13H10O3
Thermal degradation
Waste polycarbonate will degrade at high temperatures to form solid, liquid and gaseous pollutants. A study showed that the products were about 40–50 wt.% liquid, 14–16 wt.% gases, while 34–43 wt.% remained as solid residue. Liquid products contained mainly phenol derivatives (~75wt.%) and bisphenol (~10wt.%) also present. Polycarbonate, however, can be safely used as a carbon source in the steel-making industry.
Phenol derivatives are environmental pollutants, classified as volatile organic compounds (VOC). Studies show they are likely to facilitate ground level ozone formation and increase photo-chemical smog. In aquatic bodies, they can potentially accumulate in organisms. They are persistent in landfills, do not readily evaporate and would remain in the atmosphere.
Effect of fungi
In 2001 a species of fungus in Belize, Geotrichum candidum, was found to consume the polycarbonate found in compact discs (CD). This has prospects for bioremediation. However, this effect has not been reproduced.
| Physical sciences | Polymers | Chemistry |
536314 | https://en.wikipedia.org/wiki/Cinnamaldehyde | Cinnamaldehyde | Cinnamaldehyde is an organic compound with the formula C9H8O or C₆H₅CH=CHCHO. Occurring naturally as predominantly the trans (E) isomer, it gives cinnamon its flavor and odor. It is a phenylpropanoid that is naturally synthesized by the shikimate pathway. This pale yellow, viscous liquid occurs in the bark of cinnamon trees and other species of the genus Cinnamomum. It is an essential oil. The bark of cinnamon tree contains high concentrations of cinnamaldehyde.
Structure and synthesis
Cinnamaldehyde was isolated from cinnamon essential oil in 1834 by Jean-Baptiste Dumas and Eugène-Melchior Péligot and synthesized in the laboratory by the Italian chemist Luigi Chiozza in 1854.
The natural product is trans-cinnamaldehyde. The molecule consists of a benzene ring attached to an unsaturated aldehyde. Cinnamaldehyde is an α,β-unsaturated carbonyl compound. Its color is due to the π → π* transition: increased conjugation in comparison with acrolein shifts this band towards the visible.
Biosynthesis
Cinnamaldehyde is biosynthesized from phenylalanine. Deamination of L-phenylalanine into cinnamic acid is catalyzed by phenylalanine ammonia lyase (PAL). PAL catalyzes this reaction by a non-oxidative deamination. This deamination relies on the MIO prosthetic group of PAL. PAL gives rise to trans-cinnamic acid. In the second step, 4-coumarate–CoA ligase (4CL) converts cinnamic acid to cinnamoyl-CoA by an acid–thiol ligation. 4CL uses ATP to catalyze the formation of cinnamoyl-CoA. 4CL effects this reaction in two steps. 4CL forms a hydroxycinnamate–AMP anhydride, followed by a nucleophile attack on the carbonyl of the acyl adenylate. Finally, Cinnamoyl-CoA is reduced by NADPH catalyzed by CCR (cinnamoyl-CoA reductase) to form cinnamaldehyde.
Preparation
Several methods of laboratory synthesis exist. The compound can be prepared from related compounds such as cinnamyl alcohol. An early synthesis involved the aldol condensation of benzaldehyde and acetaldehyde. Cinnamaldehyde can also be obtained from the steam distillation of the oil of cinnamon bark.
Applications
As a flavorant
The most obvious application for cinnamaldehyde is as flavoring in chewing gum, ice cream, candy, e-liquid and beverages; use levels range from 9 to 4,900 parts per million (ppm) (that is, less than 0.5%). It is also used in some perfumes of natural, sweet, or fruity scents. Almond, apricot, butterscotch, and other aromas may partially employ the compound for their pleasant smells. Cinnamaldehyde can be used as a food adulterant; powdered beechnut husk aromatized with cinnamaldehyde can be marketed as powdered cinnamon. Some breakfast cereals contain as much as 187 ppm cinnamaldehyde.
As an agrichemical
Cinnamaldehyde has been tested as a safe and effective insecticide against mosquito larvae. A concentration of 29 ppm of cinnamaldehyde kills half of Aedes aegypti mosquito larvae in 24 hours. Trans-cinnamaldehyde works as a potent fumigant and practical repellant for adult mosquitos. It also has antibacterial and antifungal properties.
Miscellaneous uses
Cinnamaldehyde is a corrosion inhibitor for steel and other alloys. It is believed to form a protective film on the metal surface.
Derivatives
Numerous derivatives of cinnamaldehyde are commercially useful. Dihydrocinnamyl alcohol (3-phenylpropanol) occurs naturally but is produced by double hydrogenation of cinnamaldehyde. It has the fragrances of hyacinth and lilac. Cinnamyl alcohol similarly occurs naturally and has the odor of lilac but can be also produced starting from cinnamaldehyde. Dihydrocinnamaldehyde is produced by the selective hydrogenation of the alkene subunit. α-Amylcinnamaldehyde and α-hexylcinnamaldehyde are important commercial fragrances, but they are not prepared from cinnamaldehyde. Hydrogenation of cinnamaldehyde, if directed to the alkene, gives hydrocinnamaldehyde.
Toxicology
Cinnamaldehyde is used in agriculture because of its low toxicity, but it is a skin irritant. Cinnamaldehyde may cause allergic contact stomatitis in sensitised individuals, however allergy to the compound is believed to be uncommon.
Cinnamaldehyde can contain traces of styrene, which arises during storage or transport. Styrene especially forms in high humidity and high temperatures.
DNA repair
Cinnamaldehyde is a dietary antimutagen that effectively inhibits both induced and spontaneous mutations. Experimental evidence indicates that cinnamaldehyde induces a type of DNA damage in the bacterium Escherichia coli and in human cells that elicits recombinational DNA repair that then reduces spontaneous mutations.
In mice, X-ray–induced chromosome aberrations were reduced when cinnamaldehyde was given orally to the mice after X-ray irradiation, perhaps due to cinnamaldehyde-stimulated DNA repair.
| Physical sciences | Phenylpropanoids | Chemistry |
536505 | https://en.wikipedia.org/wiki/Rod%20%28unit%29 | Rod (unit) | The rod, perch, or pole (sometimes also lug) is a surveyor's tool and unit of length of various historical definitions. In British imperial and US customary units, it is defined as feet, equal to exactly of a mile, or yards (a quarter of a surveyor's chain), and is exactly 5.0292 meters. The rod is useful as a unit of length because integer multiples of it can form one acre of square measure (area). The 'perfect acre' is a rectangular area of 43,560 square feet, bounded by sides 660 feet (a furlong) long and 66 feet (a chain) wide (220 yards by 22 yards) or, equivalently, 40 rods by 4 rods. An acre is therefore 160 square rods or 10 square chains.
The name perch derives from the Ancient Roman unit, the pertica.
The measure also has a relationship with the military pike of about the same size. Both measures date from the sixteenth century, when the pike was still utilized in national armies. The tool has been supplanted, first by steel tapes and later by electronic tools such as surveyor lasers and optical target devices for surveying lands. In dialectal English, the term lug has also been used, although the Oxford English Dictionary states that this unit, while usually of feet, may also be of 15, 18, 20, or 21 feet.
In the United States until 1 January 2023, the rod was often defined as 16.5 US survey feet, or approximately 5.029 210 058 m.
History
In England, the perch was officially discouraged in favour of the rod as early as the 15th century; however, local customs maintained its use. In the 13th century, perches were variously recorded in lengths of , , and ; and even as late as 1820, a House of Commons report notes lengths of , , , , and even . In Ireland, a perch was standardized at , making an Irish chain, furlong and mile proportionately longer by 27.27% than the "standard" English measure.
Until English King Henry VIII seized the lands of the Roman Catholic Church in 1536, land measures as we now know them were essentially unknown. Instead a narrative system of landmarks and lists was used. Henry wanted to raise even more funds for his wars than he'd seized directly from church property (he'd also assumed the debts of the monasteries), and as James Burke writes and quotes in the book Connections that the English monk Richard Benese "produced a book on how to survey land using the simple tools of the time, a rod with cord carrying knots at certain intervals, waxed and resined against wet weather." Benese poetically described the measure of an acre in terms of a perch:
The practice of using surveyor's chains, and perch-length rods made into a detachable stiff chain, came about a century later when iron was a more plentiful and common material. A chain is a larger unit of length measuring , or 22 yards, or 100 links, or 4 rods (20.1168 meters). There are 10 chains or 40 rods in a furlong (eighth-mile), and so 80 chains or 320 rods in one statute mile (1760 yards, 1609.344 m, 1.609344 km); the definition of which was legally set in 1593 and popularized by Royal surveyor (called the 'sworn viewer') John Ogilby only after the Great Fire of London (1666).
An acre is defined as the area of 10 square chains (that is, an area of one chain by one furlong), and derives from the shapes of new-tech plows and the desire to quickly survey seized church lands into a quantity of squares for quick sales by Henry VIII's agents; buyers simply wanted to know what they were buying whereas Henry was raising cash for wars against Scotland and France. Consequently, the surveyor's chain and surveyor rods or poles (the perch) have been used for several centuries in Britain and in many other countries influenced by British practices such as North America and Australia. By the time of the industrial revolution and the quickening of land sales, canal and railway surveys, et al. Surveyor rods such as used by George Washington were generally made of dimensionally stable metal—semi-flexible drawn wrought iron linkable bar stock (not steel), such that the four folded elements of a chain were easily transportable through brush and branches when carried by a single man of a surveyor's crew. With a direct ratio to the length of a surveyor's chain and the sides of both an acre and a square (mile), they were common tools used by surveyors, if only to lay out a known plottable baseline in rough terrain thereafter serving as the reference line for instrumental (theodolite) triangulations.
The rod as a survey measure was standardized by Edmund Gunter in England in 1607 as a quarter of a chain (of ), or long.
In ancient cultures
The perch (pertica) as a lineal measure in Rome (also decempeda) was 10 Roman feet (2.96 metres), and in France varied from 10 feet (perche romanie) to 22 feet (perche d'arpent—apparently of "the range of an arrow"—about 220 feet). To confuse matters further, by ancient Roman definition, an arpent equalled 120 Roman feet. The related unit of square measure was the scrupulum or decempeda quadrata, equivalent to about .
In continental Europe
Units comparable to the perch, pole or rod were used in many European countries, with names that include and canne, , and pertica, and . They were subdivided in many different ways, and were of many different lengths.
In Britain and Ireland
In England, the rod or perch was first defined in law by the Composition of Yards and Perches, one of the statutes of uncertain date from the late 13th to early 14th centuries: tres pedes faciunt ulnam, quinque ulne & dimidia faciunt perticam (three feet make a yard, five and a half yards make a perch).
The length of the chain was standardized in 1620 by Edmund Gunter at exactly four rods. Fields were measured in acres, which were one chain (four rods) by one furlong (in the United Kingdom, ten chains).
Bars of metal one rod long were used as standards of length when surveying land. The rod was still in use as a common unit of measurement in the mid-19th century, when Henry David Thoreau used it frequently when describing distances in his work, Walden.
In traditional Scottish units, a Scottish rood (ruid in Lowland Scots, ròd in Scottish Gaelic), also fall measures 222 inches (6 ells).
Modern use
The rod was phased out as a legal unit of measurement in the United Kingdom as part of a ten-year metrication process that began on 24 May 1965.
In the United States, the rod, along with the chain, furlong, and statute mile (as well as the survey inch and survey foot) were based on the pre-1959 values for United States customary units of linear measurement until 1 January 2023. The Mendenhall Order of 1893 defined the yard as exactly meters, with all other units of linear measurement, including the rod, based on the yard. In 1959, an international agreement (the international yard and pound agreement), defined the yard as the fundamental unit of length in the Imperial/USCU system, defined as exactly 0.9144 metres. However, the above-noted units, when used in surveying, may retain their pre-1959 values, depending on the legislation in each state. The U.S. National Geodetic Survey and National Institute of Standards and Technology have replaced the definition for the above-mentioned units by the international 1959 definition of the foot, being exactly 0.3048 meters.
Despite no longer being in widespread use, the rod is still employed in certain specialized fields. In recreational canoeing, maps measure portages (overland paths where canoes must be carried) in rods; typical canoes are approximately one rod long. The term is also in widespread use in the acquisition of pipeline easements, as the offers for an easement are often expressed on a "price per rod".
In the United Kingdom, the sizes of allotment gardens continue to be measured in square poles in some areas, sometimes being referred to simply as poles rather than square poles.
In Vermont, the default right-of-way width of state and town highways and trails is three rods . Rods can also be found on the older legal descriptions of tracts of land in the United States, following the "metes and bounds" method of land survey; as shown in this actual legal description of rural real estate:
Area and volume
The terms pole, perch, rod and rood have been used as units of area, and perch is also used as a unit of volume. As a unit of area, a square perch (the perch being standardized to equal feet, or yards) is equal to a square rod, or acre. There are 40 square perches to a rood (for example a rectangular area of 40 rods times one rod), and 160 square perches to an acre (for example a rectangular area of 40 rods times 4 rods). This unit is usually referred to as a perch or pole even though square perch and square pole were the more precise terms. Rod was also sometimes used as a unit of area to refer to a rood.
However, in the traditional French-based system in some countries, 1 square perche is 42.21 square metres.
As of August 2013, perches and roods are used as government survey units in Jamaica. They appear on most property title documents. The perch is also in extensive use in Sri Lanka, being favored even over the rood and acre in real estate listings there. Perches were informally used as a measure in Queensland real estate until the early 21st century, mostly for historical gazetted properties in older suburbs.
Volume
A traditional unit of volume for stone and other masonry. A perch of masonry is the volume of a stone wall one perch () long, high, and thick. This is equivalent to exactly .
There are two different measurements for a perch depending on the type of masonry that is being built:
A dressed stone work is measured by the -cubic foot perch () long, high, and thick. This is equivalent to exactly .
a brick work or rubble wall made of broken stone of irregular size, shape and texture, made of undressed stone, is measured by the () long, high, and thick. This is equivalent to exactly .
| Physical sciences | English | Basics and measurement |
536889 | https://en.wikipedia.org/wiki/Marmoset | Marmoset | The marmosets (), also known as zaris or sagoin, are twenty-two New World monkey species of the genera Callithrix, Cebuella, Callibella, and Mico. All four genera are part of the biological family Callitrichidae. The term "marmoset" is also used in reference to Goeldi's marmoset, Callimico goeldii, which is closely related.
Most marmosets are about long. Relative to other monkeys, they show some apparently primitive features; they have claws rather than nails, and tactile hairs on their wrists. They lack wisdom teeth, and their brain layout seems to be relatively primitive. Their body temperature is unusually variable, changing by up to in a day. Marmosets are native to South America and have been found in Bolivia, Brazil, Colombia, Ecuador, Paraguay, and Peru. They have also been occasionally spotted in Central America and southern Mexico. They are sometimes kept as pets, though they have specific dietary and habitat needs that require consideration.
According to recent research, marmosets exhibit germline chimerism, which is not known to occur in nature in any primates other than callitrichids. 95% of marmoset fraternal twins trade blood through chorionic fusions, making them hematopoietic chimeras.
Etymology
Callithrix comes from Ancient Greek and means "beautiful fur".
Species list
Genus Callithrix—Atlantic marmosets
Common marmoset, Callithrix jacchus
Black-tufted marmoset, Callithrix penicillata
Wied's marmoset, Callithrix kuhlii
White-headed marmoset, Callithrix geoffroyi
Buffy-headed marmoset, Callithrix flaviceps
Buffy-tufted marmoset, Callithrix aurita
Genus Mico—Amazonian marmosets
Rio Acari marmoset, Mico acariensis
Silvery marmoset, Mico argentatus
White marmoset, Mico leucippe
Emilia's marmoset, Mico emiliae
Black-headed marmoset, Mico nigriceps
Marca's marmoset, Mico marcai
Black-tailed marmoset, Mico melanura
Santarem marmoset, Mico humeralifer
Maués marmoset, Mico mauesi
Munduruku marmoset, Mico munduruku
Gold-and-white marmoset, Mico chrysoleucos
Hershkovitz's marmoset, Mico intermedius
Satéré marmoset, Mico saterei
Rondon's marmoset, Mico rondoni
Genus Callibella—Roosmalens' dwarf marmoset
Roosmalens' dwarf marmoset, Callibella humilis
Genus Cebuella—Pygmy Marmoset
Pygmy marmoset, Cebuella pygmaea
Behavior
Marmosets are highly active, living in the upper canopy of forest trees, and feeding on insects, fruit, leaves, tack, sap, and gum. They have long lower incisors, which allow them to chew holes in tree trunks and branches to harvest the gum inside; some species are specialised feeders on gum.
Marmosets live in family groups of three to 15, consisting of one or two breeding females, an unrelated male, their offspring, and occasionally extended family members and unrelated individuals. Their mating systems are highly variable and can include monogamy, polygyny, and polyandry. In most species, fraternal twins are usually born, but triplets are not unknown. Like other callitrichines, marmosets are characterized by a high degree of cooperative care of the young and some food sharing and tolerated theft. Adult males, females other than the mother, and older offspring participate in carrying infants. Father marmosets are an exceptionally attentive example of fathers within the animal kingdom, going as far as assisting their mates in giving birth, cleaning up afterbirth, and even biting the umbilical cords attaching their newborn offspring to their mothers. Most groups scent mark and defend the edges of their ranges, but whether they are truly territorial is unclear, as group home ranges greatly overlap.
The favorite food of marmosets is carbohydrate-rich tree sap, which they reach by gnawing holes in trunks. Their territories are centered on the trees that they regularly exploit in this way. The smaller marmosets venture into the very top of forest canopies to hunt insects that are abundant there.
Marmosets use chirps, trills, and "phee" calls to communicate with each other. "Phee" calls are long-distance vocalizations that help monkeys identify each other's locations. Marmosets have been observed to use distinctive "phee" calls for the different individuals in their group, similar to a human name.
| Biology and health sciences | New World monkeys | Animals |
537019 | https://en.wikipedia.org/wiki/Temperate%20rainforest | Temperate rainforest | Temperate rainforests are rainforests with coniferous or broadleaf forests that occur in the temperate zone and receive heavy rain.
Temperate rainforests occur in oceanic moist regions around the world: the Pacific temperate rainforests of North American Pacific Northwest as well as the Appalachian temperate rainforest in the Appalachian region of the United States; the Valdivian temperate rainforests of southwestern South America; the rainforests of New Zealand and southeastern Australia; northwest Europe (small pockets in Great Britain and larger areas in Ireland, southern Norway, northern Iberia and Brittany); southern Japan; the Black Sea–Caspian Sea region from the southeasternmost coastal zone of the Bulgarian coast, through Turkey, to Georgia, and northern Iran.
The moist conditions of temperate rainforests generally support an understory of mosses, ferns and some shrubs and berries. Temperate rainforests can be temperate coniferous forests or temperate broadleaf and mixed forests.
Definition
For temperate rainforests of North America, Alaback's definition is widely recognized:
Annual precipitation over (KJ)
Mean annual temperature is between .
However, required annual precipitation depends on factors such as distribution of rain over the year, temperatures over the year and fog presence, and definitions in other regions of the world differ considerably. For example, Australian definitions are ecological-structural rather than climatic:
Closed canopy of trees excludes at least 69% of the sky.
Forest is composed mainly of tree species which do not require fire for regeneration, but with seedlings able to regenerate under shade and in natural openings.
Australian definitions would exclude some temperate rainforests of western North America that are Coast Douglas-fir dominant, such as parts of the Klamath Mountains in southern Oregon and northern California, the Puget Lowlands of western Washington and the Georgia Depression in British Columbia, as their dominant tree species, the Coast Douglas-fir, requires stand-destroying disturbance to initiate a new cohort of seedlings. The North American definition would in turn exclude a part of temperate rainforests under definitions used elsewhere.
Canopy level
For forests, canopy refers to the upper layer or habitat zone, formed by mature tree crowns and including other biological organisms (epiphytes, lianas, arboreal animals, etc.). The canopy level is the third level of the temperate rainforest. The trees forming the canopy, conifers, can stand as tall as 100 metres or more. A variety of species survive in the canopy. The tops of these trees collect most of the rain, moisture, and photosynthesis that the rainforest takes in. They form a canopy over the forest, covering about 95% of the floor during the summer.
The canopy's coverage affects the shade tolerance levels of forest floor plants. When the canopy is in full bloom, covering about 95% of the floor, plant survival decreases. Some plant species have become shade tolerant in order to survive.
The treetops take in the heavy amount of rain and keep the lower levels of the forest damp.
The canopy survives through photosynthesis. The leaves provide energy and nutrients for the trees, which provide homes and food for the forest. Through satellite data, the radiation use efficiency (RUE) calculates the annual amount of photosynthesis that occurs in temperate rainforests. A diverse amount of photosynthesis occurs based on the location and microclimates of the forest.
Distribution
North America
Pacific temperate rainforests
A portion of the temperate rain forest region of North America, the largest area of temperate zone rainforests on the planet, is the Pacific temperate rain forests ecoregion, which occur on west-facing coastal mountains along the Pacific coast of North America, from Kodiak Island in Alaska to northern California, and are part of the Nearctic realm. In the different system established by the Commission for Environmental Cooperation, this same general region is classed as the Pacific Maritime Ecozone by Environment Canada and as the Marine West Coast Forest Level II ecoregion by the United States Environmental Protection Agency. In terms of the floristic province system used by botany, the bulk of the region is the Rocky Mountain Floristic Region but a small southern portion is part of the California Floristic Province.
Sub-ecoregions of the Pacific temperate rainforest ecoregion as defined by the WWF include the Northern Pacific coastal forests, Haida Gwaii ecoregion, Vancouver Island ecoregion, British Columbia mainland coastal forests, Central Pacific coastal forests, Cascades forests, Klamath-Siskiyou coastal forests, and Northern California coastal forests ecoregions. They vary in their species composition, but are all predominantly coniferous, sometimes with an understory of broadleaved trees and shrubs. Most of the precipitation occurs in winter, similar to Mediterranean climates, but in summer, fog moisture is extracted by the trees and produces a fog drip keeping the forest moist. The Northern California coastal forests are home to the Coast Redwood (Sequoia sempervirens), the world's tallest tree. In the other ecoregions, Coast Douglas-fir (Pseudotsuga menziesii var. menziesii), Sitka Spruce (Picea sitchensis), Western Hemlock (Tsuga heterophylla) and Western redcedar (Thuja plicata) are the most important tree species. A common feature of Pacific temperate rainforests of North America is the Nurse log, a fallen tree which as it decays, provides ecological facilitation to seedlings. Trees such as the Coast Douglas-fir, Western Hemlock, Western Red Cedar, Pacific Yew, and Vine Maple are more closely related to coniferous and deciduous trees in the temperate forests of East Asia.
Some of the largest expanses of old growth are found in Olympic National Park, Mount Rainier National Park, Mount Hood National Forest, Crater Lake National Park, Tongass National Forest, Mount St. Helens National Monument, Redwood National Park, and throughout British Columbia (including British Columbia's Coastal Mountain Ranges), with the coastal Great Bear Rainforest containing the largest expanses of old growth temperate rainforest found in the world.
British Columbia's Rocky Mountains, Cariboo Mountains, Rocky Mountain Trench (east of Prince George) and the Columbia Mountains of Southeastern British Columbia (west of the Canadian Rocky Mountains that extend into parts of Idaho and Northwestern Montana in the US), which include the Selkirk Mountains, Monashee Mountains, and the Purcell Mountains, have the largest stretch of interior temperate coniferous rainforests. These inland rainforests have more continental climate with a large proportion of the precipitation falling as snow. Being closer to the Rocky Mountains, there is more of a diverse mammalian fauna. Some of the best interior rainforests are found in Mount Revelstoke National Park and Glacier National Park (Canada) in the Columbia Mountains.
Appalachian temperate rainforests
Temperate rainforests are located in the southern Appalachian Mountains where orographic precipitation increases precipitation of weather systems coming from the west and from the Gulf of Mexico. Temperate rainforest extends through the Appalachian areas of western North Carolina, southeastern Kentucky, southwest Virginia, eastern Tennessee, northern South Carolina, and northern Georgia.
Red spruce and Fraser fir are dominant canopy trees in high mountain areas. In higher elevation (over ), Fraser fir is dominant, in middle elevation () red spruce and Fraser fir grow together, and in lower elevation () red spruce is dominant. Yellow birch, mountain ash, and mountain maple grow in the understory. Younger spruce and fir and shrubs like raspberry, blackberry, hobblebush, southern mountain cranberries, red elderberry, minniebush, southern bush honeysuckle are understory vegetation. Below the spruce-fir forest, at around , are forests of American beech, yellow birch, maple birch, and oak. Skunk cabbage and ground juniper are northern species that were pushed into the areas from the north.
The mild and wet environment supports the high diversity of fungi. Over 2,000 species live in this area and scientists estimate many unidentified fungi may be there.[10]
South America
Valdivian and Magellanic temperate rainforests
The temperate rainforests of South America are located on the Pacific coast of southern Chile, on the west-facing slopes of the southern Chilean coast range, and the Andes Mountains in both Chile and Western Argentina down to the southern tip of South America, and are part of the Neotropical realm. Temperate rainforests occur in the Valdivian temperate rain forests and Magellanic subpolar forests ecoregions. The Valdivian rainforests are home to a variety of broadleaf evergreen trees, like Aextoxicon punctatum, Eucryphia cordifolia, and southern beech (Nothofagus), but include many conifers as well, notably Alerce (Fitzroya cupressoides), one of the largest tree species of the world.
The Valdivian and Magellanic temperate rainforests are the only temperate rainforests in South America. Together they are the second largest in the world, after the Pacific temperate rainforests of North America. The Valdivian forests are a refuge for the Antarctic flora, and share many plant families and genera with the temperate rainforests of New Zealand, Tasmania, and Australia. Fully half the species of woody plants are endemic to this ecoregion.
In the Valdivian region the Andean Cordillera intercepts moist westerly winds along the Pacific coast during winter and summer months; these winds cool as they ascend the mountains, creating heavy rainfall on the mountains' west-facing slopes. The northward-flowing oceanic Humboldt Current creates humid and foggy conditions near the coast. The tree line is at about 2,400 m in the northern part of the ecoregion (35°S), and descends to 1,000 m in the south of the Valdivian region. In the summer the temperature can climb to , while during winter the temperature can drop below .
Africa
Knysna-Amatole coastal rainforests (South Africa)
The temperate rainforests of South Africa are part of the Knysna-Amatole forests that are located along South Africa's Garden Route between Cape Town and Port Elizabeth on the south-facing slopes of South Africa's Drakensberg Mountains facing the Indian Ocean. There are several coniferous podocarps that grow here. This forest receives a lot of moisture as fog from the Indian Ocean, and resembles not only other temperate rainforests worldwide, but also the montane evergreen Afromontane forests that occur at higher elevations in southern and eastern Africa. A fine example of this forest is in South Africa's Tsitsikamma National Park.
Macaronesia
Azores
The rainforests of the Azores (also known as cloud forests, due to the constant cloud coverage caused by orographic lift) are found in the more humid, montane areas that transition from the lower altitude laurissilva. They are generally found at altitudes ranging from , and receive of average annual rainfall.
Despite being located in the temperate zone, the Azores rainforest is similar in many ways to the cloud forest environments of the tropics and subtropics. These pluvial montane forests hold the highest biodiversity and degree of endemism of the whole archipelago. They are dominated by dense formations of endemic juniper, laurel, holly and tree heaths with several species of epiphytic ferns and an abundance of mosses and rainforest lichens (such as Erioderma).
The climate in the rainforest is mild and cool, averaging with a narrow diurnal temperature range and temperatures that only drop below freezing in exceptional years.
Since human settlement in the 15th century, these rainforests, which once covered most of the high altitudes of the archipelago, have gradually been reduced to relics and are now found almost exclusively on three of the nine islands (Flores, Pico and Terceira). Their main threat is the expansion of cattle grazing pastures.
Europe
Temperate rainforest occurs in fragments across the north and west of Europe in countries such as southern Norway (see Scandinavian coastal conifer forests) and northern Spain. Other temperate rainforest regions include areas of south eastern Europe such as mountains on the east coast of the Adriatic Sea, surrounding North Western Bulgaria along with the Black Sea.
Atlantic Oakwood forest (Britain and Ireland)
The woodlands are variously referred to in Britain as Upland Oakwoods, Atlantic Oakwoods, Western Oakwoods or Temperate Rainforest, Caledonian forest, and colloquially as 'Celtic Rainforests'. They are also listed in the British National Vegetation Classification as British NVC community W11 and British NVC community W17 depending on the ground flora. The majority of surviving fragments of Atlantic Oakwoods in Britain occur on steep-sided slopes above rivers and lakes which have avoided clearance and intensive grazing pressure. There are notable examples on the islands and shores of Loch Maree, Loch Sunart, Loch Lomond and one of the best preserved sites on the remote Taynish Peninsula in Argyll. There are also small areas on steep-sided riverine gorges in Snowdonia and Mid Wales, such as found at the Dolmelynllyn Estate in Gwynedd.
In England, they occur in the Lake District (Borrowdale Woods) and steep sided riverine and estuarine valleys in Devon and Cornwall and the Microclimate disused slate & granite quarries in these counties. This includes the Fowey valley in Cornwall and the valley of the river Dart which flows off Dartmoor and has rainfall in excess of 2 metres per year.
Derrycunnihy Wood, located in the Killarney National Park, is the best example of the ancient damp-climate oceanic forest that covered an estimated 80 percent of Ireland prior to the arrival of humans in 7,000 BCE.
Guy Shrubsole's Lost Rainforests of Britain attempts to find, map, photograph, and restore them.
Colchian (Colchis) rainforests (Bulgaria, Turkey and Georgia)
The Colchian rainforests are found around both the southeast and west corners of the Black Sea starting in Bulgaria all the way to Turkey and Georgia and are part of the Euxine-Colchic deciduous forests ecoregion, together with the drier Euxine forests further west. The Colchian rainforests are mixed, with deciduous black alder (Alnus glutinosa), hornbeam (Carpinus betulus and C. orientalis), Oriental beech (Fagus orientalis), and sweet chestnut (Castanea sativa) together with evergreen Nordmann fir (Abies nordmanniana, the tallest tree in Europe at 78 m), Caucasian spruce (Picea orientalis) and Scots pine (Pinus sylvestris).
The refugium is the largest throughout the Western Asian/ near Eastern region. The area has multiple representatives of disjunct relict groups of plants with the closest relatives in Eastern Asia, southern Europe, and even North America. Over 70 species of forest snails of the region are endemic. Some relict species of vertebrates are Caucasian parsley frog, Caucasian salamander, Robert's snow vole and Caucasian grouse; they are almost entirely endemic groups of animals such as lizards of genus Darevskia. In general, species composition of this refugium is quite distinct and differs from that of the other Western Eurasian refugia.
Genetic data suggest that the Colchis temperate rainforest, during the Ice Age, was fragmented into smaller parts; in particular, evolutionary lineages of the Caucasian Salamander from the central and south-western Colchis remained isolated from one another during the entire Ice Age.
Fragas do Eume (Spain)
The is a natural park situated in Galicia, north-western Spain. is a Galician word for 'natural woodland', (old-growth forest) and the park is an example of a temperate rainforest in which oak (Quercus robur and Quercus pyrenaica) is the climax vegetation. The protected area extends along the valley of the river Eume within the municipalities of , , , and . Some 500 people reside within the park. The monastery of Monastery of also lies within the park.
The area was declared a natural park (a level of protection lower than national park) in 1997. It is one of six natural parks in Galicia. The European Union has recognised the park as a Site of Community Importance. There are a number of species of ferns. Invertebrate species include the Kerry slug and it is an important site for amphibians.
Vinatovača rainforest (Serbia)
The Vinatovača rainforest, alternatively spelled vintovača, is the only rainforest in Serbia. It has been left undisturbed for centuries due to strict conservation laws starting in the 17th century.
Vinatovača is situated in the central Kučaj mountains in the Upper Resava region, at an altitude between and . It is isolated and hard to reach which helped its preservation. It is believed that trees have not been cut in Vinatovača since about 1650. Being under strict protection means not only that the trees that die of old age are not being cleared or removed, but even picking herbs or mushrooms is forbidden. It is considered as an example of what central and eastern Serbia's natural look is. Beech trees are up to tall and some specimens are estimated to be over 300 years old.
Asia
Caspian Hyrcanian forest (Iran and Azerbaijan)
The Caspian Hyrcanian mixed forests ecoregion in northern Iran contains a jungle in the form of a rainforest which stretches from the east in the Khorasan province to the west in the Ardabil Province, covering the other provinces of Gilan, Mazandaran, and Golestan. The Elburz or Alborz mountain range is the highest mountain range in the Middle East which captures the moisture of the Caspian Sea to its north and forms subtropical and temperate rainforests in the northern part of Iran. The Iranians call this forest and region Shomal which means north in Persian. This forest was known for most of the history for being home to the now extinct Caspian Tiger.
In southeast Azerbaijan, this ecoregion includes the Lankaran Lowland and the Talysh Mountains, the latter being evenly divided with Iran to the south. They are deciduous forests containing tree species such as black alder (Alnus glutinosa subsp. barbata), hornbeam (Carpinus betulus and C. orientalis), Caucasian wingnut (Pterocarya fraxinifolia), chestnut-leaved oak (Quercus castaneifolia), Caucasian oak (Quercus macranthera), oriental beech (Fagus orientalis), Persian ironwood (Parrotia persica) and Persian silk tree (Albizia julibrissin).
The existing protected areas in Azerbaijan include:
Gizil-Agach State Reserve –
Hirkan National Park –
Zuvand National Park –
Girkan State Reserve –
High elevation mountain rainforests (Taiwan)
These forests are found in eastern Taiwan and Taiwan's Central Mountain Ranges, part of the Taiwan subtropical evergreen forest region covering the higher elevations. Most of the lower elevations are covered by subtropical broadleaf evergreen forests, dominated by Chinese Cryptocarya (Cryptocarya chinensis), Castanopsis hystrix and Japanese Blue Oak (Quercus glauca). Higher elevations give way to temperate forests with large stands of old growth Taiwan Cypress (Chamaecyparis taiwanensis), Camphor tree (Cinnamomum camphora), maple (Acer spp.), Chinese yew (Taxus chinensis), Taiwan Hemlock (Tsuga chinensis), and Taiwan Douglas-fir (Pseudotsuga sinensis var. wilsoniana). These higher elevation forests include also giant conifers Formosan Cypress (Chamaecyparis formosensis) and Taiwania (Taiwania cryptomerioides) Some fine examples of forests are found in Yushan (Jade Mountain) National Park and Alishan.
Baekdu Mountain Range (Taebaek and Sobaek Mountain Ranges) and South Sea forests (Korea)
The forests that cover the mountains and valleys of the Baekdu Mountain Range – from Mt. Baekdu, in the north, to Mt. Jiri, in the southwest, forming the spine of the Korean Peninsula – and the southern coast and islands of the peninsula – including Jeju Island – feature a wide variety of conifers and broadleaf trees. Much of these forests are protected in mountain and marine national forests, such as in Hallyeohaesang National Park, which encompasses of mountainous forests spread out over 69 uninhabited islands and 30 inhabited islands in Korea's South Sea that provide a home to 1,142 plant species, including major species such as red pine, black pine, common camellia, serrata oak, and cork oak, as well as rare species such as nadopungnan (sedirea japonica), daeheongnan (cymbidium nipponicum) and the Korean winter hazel. Major animals species, such as otters, small-eared cats, and badgers also call Hallyeohaesang National Park home, and overall there are 25 mammal species, 115 bird species, 16 reptile species, 1,566 insect species, and 24 freshwater fish species found among the forested, mountains islands.
Seoraksan National Park covers of mountainous forests near the eastern coast of the Korean Peninsula, and is a UNESCO designated Biosphere Preservation District. Over 2,000 animal species live in Seoraksan, including the Korean goral, musk deer, and there are also more than 1,400 rare plant species, such as the edelweiss.
Taiheiyo (Pacific) rainforests (Japan)
Southwestern Japan's Taiheiyo evergreen forests region covers much of Shikoku and Kyūshū Islands, and the Southern/Pacific Ocean-facing side of Honshu ("Taiheiyo" is the Pacific Ocean, in Japanese). Here the natural forests are mainly broadleaf evergreen in lower elevations and deciduous in higher elevations. The Hydrangea hirta species is an endemic deciduous species that can be found in this area. The limit occurs at 500–1000 metres depending on latitude. The main tree species are members of beech family (Fagaceae). In lower altitudes these include evergreen oaks (Quercus spp.), Japanese Chinquapin (Castanopsis cuspidata) and Japanese Stone Oak (Lithocarpus edulis), and in higher altitudes Japanese Blue Beech (Fagus japonica) and Siebold's beech (Fagus crenata).
Some of the best preserved examples of forest are found in Kirishima-Yaku National Park on the Island of Yakushima off of Kyūshū in a very wet climate (the annual rainfall is 4,000 to 10,000 mm depending on altitude). Because of relatively infertile soils on granite, Yakushima's forests in higher elevations are dominated by a giant conifer species, Japanese Cedar (Cryptomeria japonica), rather than deciduous forests typical of the mainland. Other areas include Mount Kirishima near Kagoshima in southern Kyūshū. On Southern Honshū, there is a forest with the Nachi Falls located in Yoshino-Kumano National Park. This particular area of Honshū has been described as one of the rainiest spots in Japan.
Eastern Himalayan broadleaf forests (Bhutan, India, Nepal)
It is a temperate broadleaf forest ecoregion found in the middle elevations of the eastern Himalayas, including parts of Nepal, India, and Bhutan.
Southern Siberian rainforest
Temperate rainforests of the Russian Far East
The Russian Far East region is the eastern-most region of both Russia and the Asian continent as a whole. The Russian federal subjects of Primorsky Krai and Khabarovsk Krai are located in the southeast of this region, with Primorsky Krai sharing a land border with China and North Korea, and both federal subjects face the Pacific ocean to the eaat and share maritime borders with Japan. The Sikhote-Alin mountain range is located here and extends for about 1000 km in a northeast direction, parallel to the coast, from near the coastal city of Vladivostok.
Whilst the mountain range ascends from sea level to a maximum altitude of around 1900 metres contains a variety of different habitats, they are located in region with a temperate climate. During the last glacial maximum (or ice age), the area was not glaciated, allowing for the development of a complex ecosystem containing species with origins in Siberia’s boreal forest and Manchuria’s subtropical forests. Temperate rainforest covers most of the mountain slopes and the biogeographic region is known as the Primorye centre of plant diversity, a biogeographic meeting point of flora and fauna from temperate, subtropical and taiga climatic regions. Historically these forests ranged from the southeastern Pacific coast of Russia, through North Korea and into northern China, however vast human development, particularly in China, has limited the forest to its current range in the Russian Far East. In 2001, UNESCO recognized a 1.5 million hectare area of forest in the central part of the Sikhote-Alin mountains as a World Heritage Site in Russia, citing the area as one of the most unique and valuable areas of intact forest in the world
Although not limited to forests, more than 2500 species of vascular plants have been described in the Primorye biographic region, of which many are considered relict and endemic species. Flora of mosses and lichens are particularly diverse. About 200 species are listed in the IUCN Red List as rare and endangered. The forests fall in the transition zone between two biomes: the southern Asian hardwood forest and the northern coniferous forest. The rainforests are a mix of deciduous broadleaf and coniferous forest, with the dominant tree species becoming more coniferous at higher elevations, and more mixed forest found at lower elevations or within mountain valleys. The most common species include the Korean pine (Pinus koraiensis) and Manchurian fir (Abies holophylla) at the lowest elevations and coastlines. Jezo spruce (Picea jezoensis) and Khingan fir (Abies nephrolepis) are common species to be found from 700–1400 metres altitude. Other tree species include Mongolian oak (Quercus mongolica), silver birch (Betula platyphylla), Scots pine (Pinus sylvestris), trembling aspen (Populus tremula), Siberian dwarf pine (Pinus pumila), Erman's birch (Betula ermanii), and Dahurian larch (Larix gmelinii), a deciduous conifer common throughout, but dominant in the northernmost reaches of the forest Other characteristic flora include various ferns, lotus, (Nelumbo nucifera) and the willow Salix arbutifolia, Taxus cuspidata, Juniperus rigida, Phellodendron amurense, Kalopanax, Aralia elata, Maackia amurensis, Alnus japonica, Actinidia kolomikta, Schisandra chinensis, Celastrus orbiculatus, Thladiantha dubia, Weigela, Eleutherococcus, Flueggea suffruticosa, Deutzia, Betula schmidtii, Carpinus cordata, Acer mandshuricum, Parthenocissus tricuspidata, Vitis amurensis, and Panax ginseng and many others.
Along with the neighbouring Amur region of Russia, the temperate rainforests of the Russian Far East hold the last remaining habitats for the critically endangered Siberian tiger, Amur leopard, and Manchurian sika deer. It has been estimated that there are less the 600 tigers. and around 90 leopards left in the wild. The area also contains populations of Asiatic black bears, Kamchatka brown bears, and Mongolian grey wolves, as the Russian Far East, altogether, might probably be the only place in the world where endangered tigers, leopards, bears, and grey wolves coexist. This region also happens to be some of the last of habitat of the Blakiston’s fish owl (Bubo blakistoni); along with being the world’s largest owl, it is unique in the way that it eats fish (primarily Masu salmon) and relies on old growth forests along river banks to hunt, nest, and breed. The Siberian grouse is similar to the spruce grouse and Franklin's grouse of North America, and can be found in the dense, remote pockets of broadleaf, coniferous and deciduous forests of Far East Russia. Common ungulates include red deer, roe deer, wild boar, Manchurian moose, and musk deer.
Oceania
Australian temperate rainforests
In Australia rainforests occur near the mainland east coast and in Tasmania. There are warm-temperate and cool-temperate rainforests. They are broadleaf evergreen forests with the exception of montane rainforests of Tasmania. Eucalypt forests are not classified as rainforests although some eucalypt forest types receive high annual rainfall (to over 2000 mm in Tasmania), and in the absence of fire they may develop to rainforest. If these widespread wet sclerophyll forests were considered rainforests, the total area of rainforest in Australia would be much larger.
Warm-temperate rainforest replaces subtropical rainforest on poorer soils or with increasing altitude and latitude in New South Wales and Victoria. Cool-temperate rainforests are widespread in Tasmania (Tasmanian temperate rainforests ecoregion) and they can be found scattered from the World Heritage listed Border Ranges National Park and Lamington National Park on the NSW/Queensland border to Otway Ranges, Strzelecki Ranges, Dandenong Ranges and Tarra Bulga in Victoria. In the northern NSW they are usually dominated by Antarctic Beech (Nothofagus moorei), in the southern NSW by Pinkwood (Eucryphia moorei) and Coachwood (Ceratopetalum apetalum) and in Victoria and Tasmania by Myrtle Beech (Nothofagus cunninghamii), Southern Sassafras (Atherosperma moschatum) and Mountain Ash (Eucalyptus regnans). The montane rainforests of Tasmania are dominated by Tasmanian endemic conifers (mainly Athrotaxis spp.). They are dominated by Ferns such as Cyathea cooperi, Cyathea australis, Dicksonia antarctica, Cyathea cunninghamii and Cyathea leichhardtiana.
New Zealand temperate rainforests
The temperate rainforests of New Zealand occur on the western shore of the South Island and on the North Island. The forests are made up of coniferous podocarps and broadleaf evergreen trees. The podocarps are abundant at lower elevations, while southern beech (Nothofagus) can be found on higher slopes and in the cooler southernmost rainforests. Ecoregions include the Fiordland temperate forests and Westland temperate rainforests.
| Physical sciences | Forests | Earth science |
537048 | https://en.wikipedia.org/wiki/Equation%20solving | Equation solving | In mathematics, to solve an equation is to find its solutions, which are the values (numbers, functions, sets, etc.) that fulfill the condition stated by the equation, consisting generally of two expressions related by an equals sign. When seeking a solution, one or more variables are designated as unknowns. A solution is an assignment of values to the unknown variables that makes the equality in the equation true. In other words, a solution is a value or a collection of values (one for each unknown) such that, when substituted for the unknowns, the equation becomes an equality.
A solution of an equation is often called a root of the equation, particularly but not only for polynomial equations. The set of all solutions of an equation is its solution set.
An equation may be solved either numerically or symbolically. Solving an equation numerically means that only numbers are admitted as solutions. Solving an equation symbolically means that expressions can be used for representing the solutions.
For example, the equation is solved for the unknown by the expression , because substituting for in the equation results in , a true statement. It is also possible to take the variable to be the unknown, and then the equation is solved by . Or and can both be treated as unknowns, and then there are many solutions to the equation; a symbolic solution is , where the variable may take any value. Instantiating a symbolic solution with specific numbers gives a numerical solution; for example, gives (that is, ), and gives .
The distinction between known variables and unknown variables is generally made in the statement of the problem, by phrases such as "an equation in and ", or "solve for and ", which indicate the unknowns, here and .
However, it is common to reserve , , , ... to denote the unknowns, and to use , , , ... to denote the known variables, which are often called parameters. This is typically the case when considering polynomial equations, such as quadratic equations. However, for some problems, all variables may assume either role.
Depending on the context, solving an equation may consist to find either any solution (finding a single solution is enough), all solutions, or a solution that satisfies further properties, such as belonging to a given interval. When the task is to find the solution that is the best under some criterion, this is an optimization problem. Solving an optimization problem is generally not referred to as "equation solving", as, generally, solving methods start from a particular solution for finding a better solution, and repeating the process until finding eventually the best solution.
Overview
One general form of an equation is
where is a function, are the unknowns, and is a constant. Its solutions are the elements of the inverse image (fiber)
where is the domain of the function . The set of solutions can be the empty set (there are no solutions), a singleton (there is exactly one solution), finite, or infinite (there are infinitely many solutions).
For example, an equation such as
with unknowns and , can be put in the above form by subtracting from both sides of the equation, to obtain
In this particular case there is not just one solution, but an infinite set of solutions, which can be written using set builder notation as
One particular solution is . Two other solutions are , and . There is a unique plane in three-dimensional space which passes through the three points with these coordinates, and this plane is the set of all points whose coordinates are solutions of the equation.
Solution sets
The solution set of a given set of equations or inequalities is the set of all its solutions, a solution being a tuple of values, one for each unknown, that satisfies all the equations or inequalities.
If the solution set is empty, then there are no values of the unknowns that satisfy simultaneously all equations and inequalities.
For a simple example, consider the equation
This equation can be viewed as a Diophantine equation, that is, an equation for which only integer solutions are sought. In this case, the solution set is the empty set, since 2 is not the square of an integer. However, if one searches for real solutions, there are two solutions, and ; in other words, the solution set is .
When an equation contains several unknowns, and when one has several equations with more unknowns than equations, the solution set is often infinite. In this case, the solutions cannot be listed. For representing them, a parametrization is often useful, which consists of expressing the solutions in terms of some of the unknowns or auxiliary variables. This is always possible when all the equations are linear.
Such infinite solution sets can naturally be interpreted as geometric shapes such as lines, curves (see picture), planes, and more generally algebraic varieties or manifolds. In particular, algebraic geometry may be viewed as the study of solution sets of algebraic equations.
Methods of solution
The methods for solving equations generally depend on the type of equation, both the kind of expressions in the equation and the kind of values that may be assumed by the unknowns. The variety in types of equations is large, and so are the corresponding methods. Only a few specific types are mentioned below.
In general, given a class of equations, there may be no known systematic method (algorithm) that is guaranteed to work. This may be due to a lack of mathematical knowledge; some problems were only solved after centuries of effort. But this also reflects that, in general, no such method can exist: some problems are known to be unsolvable by an algorithm, such as Hilbert's tenth problem, which was proved unsolvable in 1970.
For several classes of equations, algorithms have been found for solving them, some of which have been implemented and incorporated in computer algebra systems, but often require no more sophisticated technology than pencil and paper. In some other cases, heuristic methods are known that are often successful but that are not guaranteed to lead to success.
Brute force, trial and error, inspired guess
If the solution set of an equation is restricted to a finite set (as is the case for equations in modular arithmetic, for example), or can be limited to a finite number of possibilities (as is the case with some Diophantine equations), the solution set can be found by brute force, that is, by testing each of the possible values (candidate solutions). It may be the case, though, that the number of possibilities to be considered, although finite, is so huge that an exhaustive search is not practically feasible; this is, in fact, a requirement for strong encryption methods.
As with all kinds of problem solving, trial and error may sometimes yield a solution, in particular where the form of the equation, or its similarity to another equation with a known solution, may lead to an "inspired guess" at the solution. If a guess, when tested, fails to be a solution, consideration of the way in which it fails may lead to a modified guess.
Elementary algebra
Equations involving linear or simple rational functions of a single real-valued unknown, say , such as
can be solved using the methods of elementary algebra.
Systems of linear equations
Smaller systems of linear equations can be solved likewise by methods of elementary algebra. For solving larger systems, algorithms are used that are based on linear algebra. See Gaussian elimination and numerical solution of linear systems.
Polynomial equations
Polynomial equations of degree up to four can be solved exactly using algebraic methods, of which the quadratic formula is the simplest example. Polynomial equations with a degree of five or higher require in general numerical methods (see below) or special functions such as Bring radicals, although some specific cases may be solvable algebraically, for example
(by using the rational root theorem), and
(by using the substitution , which simplifies this to a quadratic equation in ).
Diophantine equations
In Diophantine equations the solutions are required to be integers. In some cases a brute force approach can be used, as mentioned above. In some other cases, in particular if the equation is in one unknown, it is possible to solve the equation for rational-valued unknowns (see Rational root theorem), and then find solutions to the Diophantine equation by restricting the solution set to integer-valued solutions. For example, the polynomial equation
has as rational solutions and , and so, viewed as a Diophantine equation, it has the unique solution .
In general, however, Diophantine equations are among the most difficult equations to solve.
Inverse functions
In the simple case of a function of one variable, say, , we can solve an equation of the form for some constant by considering what is known as the inverse function of .
Given a function , the inverse function, denoted and defined as , is a function such that
Now, if we apply the inverse function to both sides of , where is a constant value in , we obtain
and we have found the solution to the equation. However, depending on the function, the inverse may be difficult to be defined, or may not be a function on all of the set (only on some subset), and have many values at some point.
If just one solution will do, instead of the full solution set, it is actually sufficient if only the functional identity
holds. For example, the projection defined by has no post-inverse, but it has a pre-inverse defined by . Indeed, the equation is solved by
Examples of inverse functions include the th root (inverse of ); the logarithm (inverse of ); the inverse trigonometric functions; and Lambert's function (inverse of ).
Factorization
If the left-hand side expression of an equation can be factorized as , the solution set of the original solution consists of the union of the solution sets of the two equations and .
For example, the equation
can be rewritten, using the identity as
which can be factorized into
The solutions are thus the solutions of the equation , and are thus the set
Numerical methods
With more complicated equations in real or complex numbers, simple methods to solve equations can fail. Often, root-finding algorithms like the Newton–Raphson method can be used to find a numerical solution to an equation, which, for some applications, can be entirely sufficient to solve some problem.
There are also numerical methods for systems of linear equations.
Matrix equations
Equations involving matrices and vectors of real numbers can often be solved by using methods from linear algebra.
Differential equations
There is a vast body of methods for solving various kinds of differential equations, both numerically and analytically. A particular class of problem that can be considered to belong here is integration, and the analytic methods for solving this kind of problems are now called symbolic integration. Solutions of differential equations can be implicit or explicit.
| Mathematics | Basics | null |
537218 | https://en.wikipedia.org/wiki/Gun%20dog | Gun dog | Gun dogs or bird dogs are types of hunting dogs developed to assist hunters in finding and retrieving game, typically various fowls that are shot down on the wing (in flight). The term hunting dog is broad and includes all breeds and skills of hunting canines, but "gun dogs" refers to canines that are trained to work alongside a loud firearm while hunting or retrieving game. Gun dogs are divided into three main categories: pointers and retrievers, setters and spaniels, and water dogs.
Types
There are several breeds used as gundogs with varying instinctive skill sets for each one. Gun dogs are divided into three primary types: retrievers, pointing breeds, and water dogs.
Method of work
Depending upon how they are trained, dogs may be useful in a variety of hunting situations.
Pointers and setters
Upon reaching the field, the handler often will cast or direct the dog in a wide circle. Experienced dogs will search the edges of the field knowing that birds are usually found there. This wide run helps to burn off the dog's initial exuberance and may help the dog establish its bearings and form a "background" upon which game smells will be processed. The dog then begins working back and forth, starting near the hunter and slowly ranging out. The dog repeats this process as the hunters move through the field. How far a handler allows the dog to range is a matter of personal preference. When a pair of dogs work as a team, one works close in while the other ranges out in larger circles. If either dog becomes birdy, the other dog works its way over to assist. Good bird dogs are alert to their handlers and to the disposition of other dogs in the field. They should readily comply if the handler casts them to an area of particular interest, such as a brush pile.
When game is detected, a dog freezes, either pointing or crouching. If other dogs are present, they also freeze, "honoring" the first dog's point. The pointing dog remains motionless until the hunters are in position. Handlers give the command whoa, instructing the dog to remain still. What happens next depends on the dog's training. Some trainers train the dog to stay motionless while the hunter steps forward and flushes the game. Other trainers direct the dog to flush the game with a command such as get it! Pointing dogs excel on covey type birds such as bobwhite, quail, and grouse as these birds will hold in position well, allowing the hunter to approach and get into position.
If a bird is downed, the dogs are instructed to search for and retrieve it with the command dead bird, or simply dead.
Water (flushing) dogs
When hunting upland game, flushing dogs (spaniels and retrievers) work much more closely with the hunter. Flushers will not cover the same amount of ground as a pointing dog as the flusher must be kept within shotgun distance. Flushing dogs are often used on birds that run from the hunter. On such birds as pheasant, an aggressive flush is necessary to spring the bird to wing. Flushing dogs excel on these types of bird because they do not point the birds, giving them little time for escape on the ground. Pointing breeds are used on such birds, but must be well trained to know when the bird pointed has moved.
Once a bird has been flushed, the dog will sit or "hup" to watch the flight of the bird and mark the fallen birds for retrieval. A dog which does this successfully is referred to as "steady to wing and shot". Steadiness is the hallmark of the finished spaniel.
When a bird is shot, the dog should mark where it fell and wait until given the command to retrieve. Once commanded, the dog will race to the point of fall, pick up the bird, and return it to the handler.
Retrievers
Retrievers are typically used when waterfowl hunting, although they can also be employed in hunting upland birds as well. Since a majority of waterfowl hunting employs the use of small boats in winter conditions, retrievers are expected to remain sitting calmly and quietly until sent to retrieve. As birds move into range, a well-trained retriever will watch and follow the handler's gun as he shoots, marking, and remembering each bird that is downed. This is called "marking off the gun", and the downed birds are called "marks". Retrievers often are expected to recall the location of many downed birds in one sitting before they retrieve game. Once the shooting has ceased, the handler commands the dog to retrieve each bird that has been downed. If a dog did not see the bird fall, a retriever takes direction from the handler, who can use hand and whistle signals to guide the dog to the unseen downed bird. This is called a "blind" retrieve. During a typical day of shooting, additional birds are frequently downed while the dog is performing a retrieve. Retrievers are taught to ignore these "diversions" until the current retrieve has been completed. Also at times multiple dogs are used on a hunt, and retrievers are also taught to "honor" another dog's retrieve by remaining calm and quiet while the other dog is working.
List of gundog breeds
Barbet
Bracco Italiano
Braque du Bourbonnais
Brittany
Burgos Pointer
Drentsche Patrijshond
Cesky Fousek
Corded Poodle
German Longhaired Pointer
German Shorthaired Pointer
German Spaniel
German Wirehaired Pointer
Hungarian Wirehaired Vizsla
Kooikerhondje
Lagotto Romagnolo
Large Münsterländer
Pointer
Portuguese Water Dog
Portuguese Pointer
Pudelpointer
Chesapeake Bay Retriever
Curly Coated Retriever
Flat-Coated Retriever
Golden Retriever
Labrador Retriever
Nova Scotia Duck Tolling Retriever
English Setter
Gordon Setter
Irish Setter
Irish Red and White Setter
Small Münsterländer
American Cocker Spaniel
American Water Spaniel
Boykin Spaniel
Clumber Spaniel
English Cocker Spaniel
English Springer Spaniel
Field Spaniel
Irish Water Spaniel
Sussex Spaniel
Welsh Springer Spaniel
Spanish Water Dog
Spinone Italiano
Stichelhaar
Standard Poodle
Vizsla
Weimaraner
Wirehaired Pointing Griffon
Wirehaired Vizsla
Kennel club classification
When competing in conformation shows most kennel clubs, including The Kennel Club group pedigree gun dog breeds together in their own gun dog group, whilst some such as the American Kennel Club group them in a sporting group.
| Biology and health sciences | Dogs | null |
20173627 | https://en.wikipedia.org/wiki/Alpine%20goat | Alpine goat | The Alpine is a medium to large sized breed of domestic goat known for its very good milking ability. They have no set colours or markings (although certain markings are discriminated against). They have horns, a straight profile and erect ears.
The breed originated in the French Alps. Mature goats weigh around 61 kg (135 lbs), and are about 76 cm (30 in) tall at the shoulder. Alpine goats can range from white or gray to brown and black. Alpine goats are heavy milkers. The milk can be made into butter, cheese, soap, ice cream or any other dairy product normally made from cow's milk. They are often used for commercial dairy production, as well as homestead milk goats.
Types
Several sub-types of Alpine goats have emerged, namely:
Purebred (French) Alpines: the original type from the French Alps
American Alpines: Alpines with other genetic influences after their introduction to the United States. American Alpines have much of the visual type and temperament as French Alpines, but may have less standard markings or conformation due to crossbreeding.
Characteristics
Alpine goats are a medium to large sized breed. Males are over 81 cm (32 in) tall at the withers and females are over 76 cm (30 in) tall at the withers. Their hair is short to medium in length, and they come in all colours and combinations of colours. They have erect ears and a straight profile, and are described as being "alertly graceful" with the ability to adapt to any climate thanks to their hardy nature. They are the only breed with erect ears that comes in all colours and combinations of colours.
The sexual maturation rate among Alpine goats is at four to five months after birth for buck kids, and five to six months after birth for doe kids. However, doe kids should not be bred until they are at least 75-80 lbs. A doe’s gestation lasts for 145 – 155 days, with 150 being the average. Twins are the most common, but they can have singles, all the way up to quintuplets .
Alpine goats are friendly and highly curious, however they can be independent and strong-willed.
The American Dairy Goat Association faults all-white and Toggenburg patterned individuals.
History
Alpines goats originated in the French Alps.
Milk
Known for its milk, the Alpine goat is famous for its rich dairy production. Alpine goats are extremely popular in the dairy industry for their docile temperament, high quality milk output and long lactation. Alpine milk has relatively low fat content, with an average fat percent of 3.4%. It is higher in sugars than cows' milk but balances itself in terms of the amount of protein. Alpine goats' milk has 2.3 g of protein per 250 ml while cow’s milk has 3.4. A higher protein count is not always good, since it packs more calories with an increased fat content. Compared to Saanen goat milk, it is higher in all nutritional aspects, except the fat content, making it a much healthier choice.
Alpine goats are one of the top milk producers, alongside Saanen and Toggenburg goats. They are distinct from the other two due to their low value of fat content. This could be a direct correlation between the weight of the animal and its habitual environment. However, the Nubian goat's weight is similar to that of the Alpines at maturity, yet it produces less milk, containing a higher proportion of fat, than the Alpine goats.
The peak periods for milk production occur after four to six weeks of parturition (kidding). The optimal weight at which a goat produces optimal milk production is at least 130 pounds. For the Alpine goat that number is higher at 135 pounds and produces 2,134 pounds of milk per lactation. Good nutrition, proper milking procedures, reproductive management, and disease control are also factors that contribute to milk production of the Alpine Goat.
There are four requirements that need to be efficient for optimal dairy production. Dairy goats must be housed in specific conditions so that their milk production is not alarmed by changes. Changes in external factors can cause a decrease in milk production due to the pressure applied on the goat to adapt to these changes. The four factors for optimal production are: adequate ventilation, dry beds, uncontaminated feeder and water supply, minimal labor and disturbance.
Alpine milk, as with all goat milk, must be filtered and chilled immediately upon separation from the lactating doe when intended for human consumption. The temperature at which milk will remain the best is at . Cooling is required immediately of the milk so that there is no excess bacteria growth. Warm bacteria grows at a faster rate and multiplies so that the milk is spoiled. The milk that is refrigerated has a shelf life of about three to four weeks. Freezing the milk increases its shelf life by about four to five weeks.
| Biology and health sciences | Goats | Animals |
20180633 | https://en.wikipedia.org/wiki/Curtain%20wall%20%28fortification%29 | Curtain wall (fortification) | A curtain wall is a defensive wall between fortified towers or bastions of a castle, fortress, or town.
Ancient fortifications
Evidence for curtain walls or a series of walls surrounding a town or fortress can be found in the historical sources from Assyria and Egypt. Some notable examples are ancient Tel Lachish in Israel and Buhen in Egypt. Curtain walls were built across Europe during the Roman Empire; the early 5th century Theodosian Walls of Constantinople influenced the builders of medieval castles many centuries later.
Curtain wall castles
In medieval castles, the area surrounded by a curtain wall, with or without towers, is known as the bailey. The outermost walls with their integrated bastions and wall towers together make up the enceinte or main defensive line enclosing the site.
In medieval designs of castle and town, the curtain walls were often built to a considerable height and were fronted by a ditch or moat to make assault difficult. Walls were topped with battlements which consisted of a parapet, which was generally crenellated with merlons to protect the defenders and lower crenels or embrasures which allowed them to shoot from behind cover; merlons were sometimes pierced by loopholes or arrowslits for better protection. Behind the parapet was a wall walk from which the defenders could fight or move from one part of the castle to another. Larger curtain walls were provided with mural passages or galleries built into the thickness of the walls and provided with arrowslits. If an enemy reached the foot of the wall, they became difficult to see or shoot at directly, so some walls were fitted with a projecting wooden platform called a hoarding or brattice. Stone machicolations performed a similar function.
Early modern fortifications
The introduction of gunpowder made tall castle walls vulnerable to fire from heavy cannon, which prompted the trace italienne style from the 16th century. In these fortifications, the height of the curtain walls was reduced, and beyond the ditch, additional outworks such as ravelins and tenailles were added to protect the curtain walls from direct cannonading.
| Technology | Fortification | null |
27499693 | https://en.wikipedia.org/wiki/Partial%20permutation | Partial permutation | In combinatorial mathematics, a partial permutation, or sequence without repetition, on a finite set S
is a bijection between two specified subsets of S. That is, it is defined by two subsets U and V of equal size, and a one-to-one mapping from U to V. Equivalently, it is a partial function on S that can be extended to a permutation.
Representation
It is common to consider the case when the set S is simply the set {1, 2, ..., n} of the first n integers. In this case, a partial permutation may be represented by a string of n symbols, some of which are distinct numbers in the range from 1 to and the remaining ones of which are a special "hole" symbol ◊. In this formulation, the domain U of the partial permutation consists of the positions in the string that do not contain a hole, and each such position is mapped to the number in that position. For instance, the string "1 ◊ 2" would represent the partial permutation that maps 1 to itself and maps 3 to 2.
The seven partial permutations on two items are
◊◊, ◊1, ◊2, 1◊, 2◊, 12, 21.
Combinatorial enumeration
The number of partial permutations on n items, for n = 0, 1, 2, ..., is given by the integer sequence
1, 2, 7, 34, 209, 1546, 13327, 130922, 1441729, 17572114, 234662231, ...
where the nth item in the sequence is given by the summation formula
in which the ith term counts the number of partial permutations with support of size i, that is, the number of partial permutations with i non-hole entries.
Alternatively, it can be computed by a recurrence relation
This is determined as follows:
partial permutations where the final elements of each set are omitted:
partial permutations where the final elements of each set map to each other.
partial permutations where the final element of the first set is included, but does not map to the final element of the second set
partial permutations where the final element of the second set is included, but does not map to the final element of the first set
, the partial permutations included in both counts 3 and 4, those permutations where the final elements of both sets are included, but do not map to each other.
Restricted partial permutations
Some authors restrict partial permutations so that either the domain
or the range of the bijection is forced to consist of the first k items in the set of n items being permuted, for some k. In the former case, a partial permutation of length k from an n-set is just a sequence of k terms from the n-set without repetition. (In elementary combinatorics, these objects are sometimes confusingly called "k-permutations" of the n-set.)
| Mathematics | Combinatorics | null |
34335892 | https://en.wikipedia.org/wiki/Climate%20change%20in%20Africa | Climate change in Africa | Climate change in Africa is an increasingly serious threat as Africa is among the most vulnerable continents to the effects of climate change. Some sources even classify Africa as "the most vulnerable continent on Earth". Climate change and climate variability will likely reduce agricultural production, food security and water security. As a result, there will be negative consequences on people's lives and sustainable development in Africa.
Over the coming decades, warming from climate change is expected across almost all the Earth's surface, and global mean rainfall will increase. Currently, Africa is warming faster than the rest of the world on average. Large portions of the continent may become uninhabitable as a result of the rapid effects of climate change, which would have disastrous effects on human health, food security, and poverty. Regional effects on rainfall in the tropics are expected to be much more spatially variable. The direction of change at any one location is often less certain.
Observed surface temperatures have generally increased by about 1 °C in Africa since the late 19th century to the early 21st century. In the Sahel, the increase has been as much as 3 °C for the minimum temperature at the end of the dry season. Data for temperature and rainfall shows discrepancies from the norm, both in timing and location.
For instance, Kenya has a high vulnerability to the impacts of climate change. The main climate hazards include droughts and floods as rainfall will likely become more intense and less predictable. Climate models predict that temperatures will rise by 0.5 to 2 °C. In the informal urban settlements of Nairobi the urban heat island effect adds to the problem as it creates even warmer ambient temperatures. This is due to home construction materials, lack of ventilation, sparse green space, and poor access to electrical power and other services.
The African Union has put forward 47 goals and corresponding actions in a 2014 draft report to combat and mitigate climate change in Africa. The International Monetary Fund suggested in 2021 that $50 billion might be necessary to cover the costs of climate change adaptation in Africa.
Greenhouse gas emissions
Africa's per person greenhouse gas emissions are low compared to other continents. Emissions from land use change are uncertain, especially in Central Africa. The main source of uncertainty comes from carbon dioxide fluxes in the LULUCF sector (this acronym stands for land use, land-use change, and forestry).
Impacts
Temperature and weather changes
Observed surface temperatures have generally increased over Africa since the late 19th century to the early 21st century by about 1 °C, but locally as much as 3 °C for minimum temperature in the Sahel at the end of the dry season. Observed precipitation trends indicate spatial and temporal discrepancies as expected. The observed changes in temperature and precipitation vary regionally.
Current climate models (as summarised in the IPCC Sixth Assessment Report) predict increases in frequency and intensity of drought and heavy rainfall events. They also predict decreases in mean precipitation almost everywhere in Africa, with medium to high confidence. However, local rainfall trends and socio-climatic interactions are likely to manifest in mixed patterns. Therefore, the converging impacts of climate change will vary across the continent. In rural areas, rainfall patterns influence water usage.
A study in 2019 predicted increased dry spell length during wet seasons and increased extreme rainfall rates in Africa. In other words: "both ends of Africa's weather extremes will get more severe". The research found that most climate models will not be able to capture the extent of these changes because they are not convection-permitting at their coarse grid scales.
Sea level rise
In Africa, future population growth amplifies risks from sea level rise. Some 54.2 million people lived in the highly exposed low elevation coastal zones (LECZ) around 2000. This number will effectively double to around 110 million people by 2030. By 2060 it will be around 185 to 230 million people, depending on the extent of population growth. The average regional sea level rise will be around 21 cm by 2060. At that point climate change scenarios will make little difference. But local geography and population trends interact to increase the exposure to hazards like 100-year floods in a complex way.
In the near term, some of the largest displacement is projected to occur in the East Africa region. At least 750,000 people there are likely to be displaced from the coasts between 2020 and 2050. Scientific studies estimate that 12 major African cities would collectively sustain cumulative damages of US$65 billion for the "moderate" climate change scenario RCP4.5 by 2050. These cities are Abidjan, Alexandria, Algiers, Cape Town, Casablanca, Dakar, Dar es Salaam, Durban, Lagos, Lomé, Luanda and Maputo. Under the high-emission scenario RCP8.5 the damage would amount to US$86.5 billion. The version of the high-emission scenario with additional impacts from high ice sheet instability would involve up to US$137.5 billion in damages. The damage from these three scenarios accounting additionally for "low-probability, high-damage events" would rise to US$187 billion, US$206 billion and US$397 billion respectively. In these estimates, the Egyptian city of Alexandria alone accounts for around half of this figure. Hundreds of thousands of people in its low-lying areas may already need relocation in the coming decade. Across sub-Saharan Africa as a whole, damage from sea level rise could reach 2–4% of GDP by 2050. However this figure depends on the extent of future economic growth and adaptation.
In the longer term, Egypt, Mozambique and Tanzania are likely to have the largest number of people affected by annual flooding amongst all African countries. This projection assumes global warming will reach 4 °C by the end of the century. That rise is associated with the RCP8.5 scenario. Under RCP8.5, 10 important cultural sites would be at risk of flooding and erosion by the end of the century. These are the Casbah of Algiers, Carthage Archaeological site, Kerkouane, Leptis Magna Archaeological site, Medina of Sousse, Medina of Tunis, Sabratha Archaeological site, Robben Island, Island of Saint-Louis and Tipasa. A total of 15 Ramsar sites and other natural heritage sites would face similar risks. These are Bao Bolong Wetland Reserve, Delta du Saloum National Park, Diawling National Park, Golfe de Boughrara, Kalissaye, Lagune de Ghar el Melh et Delta de la Mejerda, Marromeu Game Reserve, Parc Naturel des Mangroves du Fleuve Cacheu, Seal Ledges Provincial Nature Reserve, Sebkhet Halk Elmanzel et Oued Essed, Sebkhet Soliman, Réserve Naturelle d'Intérêt Communautaire de la Somone, Songor Biosphere Reserve, Tanbi Wetland Complex and Watamu Marine National Park.
Socioeconomic impacts
Climate change will increasingly impact Africa due to many factors. These impacts are already being felt and will increase in magnitude if action is not taken to reduce global carbon emissions. The impacts include higher temperatures, drought, changing rainfall patterns, and increased climate variability. These conditions have a bearing on energy production and consumption. The recent drought in many African countries, which has been linked to climate change, adversely affected both energy security and economic growth across the continent.
Africa will be one of the regions most impacted by the adverse effects of climate change. Reasons for Africa's vulnerability are diverse and include low levels of adaptive capacity, poor diffusion of technologies and information relevant to supporting adaptation, and high dependence on agro-ecosystems for livelihoods. Many countries across Africa are classified as Least-Developed Countries (LDCs) with poor socio-economic conditions, and by implication are faced with particular challenges in responding to the impacts of climate change.
Pronounced risks identified for Africa in the IPCC's Fifth Assessment Report relate to ecosystems, water availability, and agricultural systems, with implications for food security.
In 2022, over 6,000 respondents from ten African nations took part in a climate survey conducted by the European Investment Bank. The survey found that 88% of respondents claimed climate change was hurting their lives, while 61% of respondents claimed that environmental destruction has impacted their income or source of livelihood. These losses are usually the result of severe drought, increasing sea levels or coastal erosion, or extreme weather events like floods or storms.
More than half of African respondents (57%) said that they or people they know have already made steps to adapt to the effects of climate change. Among these measures are investments in water-saving devices to mitigate the effects of drought and drain clearance ahead of flooding. 34% of all African respondents said climate change is one of the most pressing issues confronting their country, among other key issues such as inflation and access to health care.
Economic impacts
Africa is warming faster than the rest of the world on average. Large portions of the continent may become uninhabitable as a result and Africa's gross domestic product (GDP) may decline by 2% as a result of a 1 °C rise in average world temperature, and by 12% as a result of a 4 °C rise in temperature. Crop yields are anticipated to drastically decrease as a result of rising temperatures and it is anticipated that heavy rains would fall more frequently and intensely throughout Africa, increasing the risk of floods.
Additionally, Africa loses between $7 billion and $15 billion a year due to climate change, projected to reach up to $50 billion by 2030.
Agriculture
Agriculture is a particularly important sector in Africa, contributing towards livelihoods and economies across the continent. On average, agriculture in Sub-Saharan Africa contributes 15% of the total GDP. Africa's geography makes it particularly vulnerable to climate change, and 70% of the population rely on rain-fed agriculture for their livelihoods. Smallholder farms account for 80% of cultivated lands in Sub-Saharan Africa. The IPCC in 2007 projected that climate variability and change would severely compromise agricultural productivity and access to food. This projection was assigned "high confidence". Cropping systems, livestock and fisheries will be at greater risk of pest and diseases as a result of future climate change. Crop pests already account for approximately 1/6th of farm productivity losses. Climate change will accelerate the prevalence of pests and diseases and increase the occurrence of highly impactful events. The impacts of climate change on agricultural production in Africa will have serious implications for food security and livelihoods. Between 2014 and 2018, Africa had the highest levels of food insecurity in the world.
In relation to agricultural systems, heavy reliance on rain-fed subsistence farming and low adoption of climate smart agricultural practices contribute to the sector's high levels of vulnerability. The situation is compounded by poor reliability of, and access to, climate data and information to support adaptation actions. Observed and projected disruptions in precipitation patterns due to climate change are likely to shorten growing seasons and affect crop yield in many parts of Africa. Furthermore, the agriculture sector in Africa is dominated by smallholder farmers with limited access to technology and the resources to adapt.
Climate variability and change have been and continue to be the principal source of fluctuations in global food production across developing countries where production is highly rain-dependent. The agriculture sector is sensitive to climate variability, especially the inter-annual variability of precipitation, temperature patterns, and extreme weather events (droughts and floods). These climatic events are predicted to increase in the future and are expected to have significant consequences to the agriculture sector. This would have a negative influence on food prices, food security, and land-use decisions. Yields from rainfed agriculture in some African countries could be reduced by up to 50% by 2020. To prevent the future destructive impact of climate variability on food production, it is crucial to adjust or suggest possible policies to cope with increased climate variability. African countries need to build a national legal framework to manage food resources in accordance with the anticipated climate variability. However, before devising a policy to cope with the impacts of climate variability, especially to the agriculture sector, it is critical to have a clear understanding of how climate variability affects different food crops. This is particularly relevant in 2020 due to the severe invasion of Locusts adversely affecting agriculture in eastern Africa. The invasion was partially attributed to climate change – the warmer temperature and heavier rainfall which caused an abnormal increase in the number of locusts.
In East Africa, climate change is anticipated to intensify the frequency and intensity of drought and flooding, which can have an adverse impact on the agricultural sector. Climate change will have varying effects on agricultural production in East Africa. Research from the International Food Policy Research Institute (IFPRI) suggest an increase in maize yields for most East Africa, but yield losses in parts of Ethiopia, Democratic Republic of Congo (DRC), Tanzania and northern Uganda. Projections of climate change are also anticipated to reduce the potential of the cultivated land to produce crops of high quantity and quality.
Climate change in Kenya is expected to have large impacts on the agricultural sector, which is predominantly rain-fed and thus highly vulnerable to changes in temperature and rainfall patterns, and extreme weather events. Impacts are likely to be particularly pronounced in the arid and semi-arid lands (ASALs) where livestock production is the key economic and livelihood activity. In the ASALs, over 70% of livestock mortality is a result of drought. Over the next 10 years, 52% of the ASAL cattle population are at risk of loss because of extreme temperature stress.
Climate change will exacerbate the vulnerability of the agricultural sector in most Southern African countries which are already limited by poor infrastructure and a lag in technological inputs and innovation. Maize accounts for nearly half of the cultivated land in Southern Africa, and under future climate change, yields could decrease by 30%. Temperatures increases also encourage a wide spread of weeds and pests.
Climate change will significantly affect agriculture in West Africa by increasing the variability in food production, access and availability.
Higher rainfall intensity, prolonged dry spells and high temperatures are expected to negatively impact cassava, maize and bean production in Central Africa. Floods and erosion occurrence are expected to damage the already limited transportation infrastructure in the region leading to post harvest losses. Exportation of economic crops like coffee and cocoa are on the rise within the region but these crops are highly vulnerable to climate change. Conflicts and political instability have had an impact on agriculture contribution to the regional GDP and this impact will be exacerbated by climatic risks.
Africa's gross domestic product (GDP) may decline by 2% as a result of a 1 °C rise in average world temperature, and by 12% as a result of a 4 °C rise in temperature. Crop yields are anticipated to drastically decrease as a result of rising temperatures and an increase in the likelihood of drought throughout the continent. Additionally, it is anticipated that heavy rains would fall more frequently and intensely throughout Africa, increasing the risk of floods.
Energy
With increasing population and corresponding energy demand, energy security must be addressed because energy is crucial for sustainable development. Climate change has affected energy sectors in Africa as many countries depend on hydropower generation. Decreasing rainfall levels and droughts have resulted in lower water levels in dams with adverse impacts on hydropower generation. This has resulted in low electrical energy production, high cost of electricity and power outages or load-shedding in some African countries that depend on hydroelectric power generation. Disruptions in hydropower generation have negatively affected various sectors in countries such as Ghana, Uganda, Kenya, and Tanzania.
Water scarcity
Water quality and availability have deteriorated in most areas of Africa, particularly due to climate change. Water resources are vulnerable and have the possibility of being strongly impacted by climate change with vast ramifications on human societies. The IPCC predicts millions of people in Africa will persistently face increased water stress due to climate variability and change (IPCC 2013). Changes in precipitation patterns directly affect surface runoff and water availability.
Climate change is likely to further exacerbate water-stressed catchments across Africa – for example the Rufiji basin in Tanzania – owing to diversity of land uses, and complex sociopolitical challenges.
Health impacts
African countries have the least efficient public health systems in the world. Infectious disease burdens such as malaria, schistosomiasis, dengue fever, meningitis, which are sensitive to climate impacts, are highest in the sub-Saharan African region. For instance, over 90 percent of annual global malaria cases are in Africa. Changes in climate will affect the spread of infectious agents as well as alter people's disposition to these infections.
According to the IPCC's Sixth Assessment Report, climate change poses a significant threat to the health of tens of millions of Africans, as it exposes them to non-optimal temperatures, extreme weather, and an increased range and transmission rate of infectious diseases.
Climate change, and resulting in increased temperatures, storms, droughts, and rising sea levels, will affect the incidence and distribution of infectious disease across the globe.
In July 2021, the World Food Programme (WFP) blamed the ongoing southern Madagascar food crisis as being caused solely by climate change and not by war or conflict. It was declared to be first famine caused by climate change.
Malaria
In Africa malaria continues to have dramatic effects on the population. As climate change continues, the specific areas likely to experience the year-round, high-risk transmission of malaria will shift from coastal West Africa to an area between the Democratic Republic of the Congo and Uganda, known as the African Highlands.
Scientific limitations when examining shifting malaria transmission rates in the African Highlands are similar to those related to broader understandings of climate change and malaria. While modeling with temperature changes shows that there is a relationship between an increase in temperature and an increase in malaria transmission, limitations still exist. Future population shifts that affect population density, as well as changes in the behavior of mosquitos, can affect transmission rates and are limiting factors in determining the future risk of malaria outbreaks, which also affect planning for correct outbreak response preparation.
With regards to malaria transmission rates in the African Highlands, factors and exposures resulting from drastic environmental changes like warmer climates, shifts in weather patterns, and increases in human impact such as deforestation, provide appropriate conditions for malaria transmission between carrier and host. Specifically, malaria is caused by the Plasmodium falciparum and Plasmodium vivax parasites which are carried by the vector Anopheles mosquito. Even though the Plasmodium vivax parasite can survive in lower temperatures, the Plasmodium falciparum parasite will only survive and replicate in the mosquito when climate temperatures are above 20 °C. Increases in humidity and rain also contribute to the replication and survival of this infectious agent. Exposure to malaria will become a greater risk to humans as the number of female Anopheles mosquitos infected with either the Plasmodium falciparum or Plasmodium vivax parasite increases.
Studies show an overall increase in climate suitability for malaria transmission resulting in an increase in the population at risk of contracting the disease. Of significant importance is the increase of epidemic potential at higher altitudes (like the African Highlands). Rising temperatures in these areas have the potential to change normally non-malarial areas to areas with seasonal epidemics. Consequently, new populations will be exposed to the disease resulting in healthy years lost. In addition, the disease burden may be more detrimental to areas that lack the ability and resources to effectively respond to such challenges and stresses.
As climate change shifts geographic areas of transmission to the African Highlands, the challenge will be to find and control the vector in areas that have not seen it before.
Impacts on conflicts and migration
The United Nations Environment Programme produced a post-conflict environmental assessment of Sudan in 2007. According to this report, environmental stresses in Sudan are interlinked with other social, economic and political issues, such as population displacement and competition over natural resources. Regional climate change, through decreased precipitation, was thought to have been one of the factors which contributed to the conflict in Darfur. Along with other environmental issues, climate change could negatively affect future development in Sudan. One of the recommendations made by UNEP was for the international community to assist Sudan in adapting to climate change.
Impacts by region
Central Africa
Central Africa, for the most part, is landlocked and is geographically threatened by climate change. Due to its high climate variability and rainfed agriculture, Central Africa is expected to experience longer and more frequent heatwaves as well as an increase in wet extremes. The global mean temperature in this region is to increase by 1.5 °C to 2 °C.
The carbon dioxide-absorbing capacity of forests in the Congo Basin have decreased. This decrease has occurred due to increasing heat and drought causing decreased tree growth. This suggests that even unlogged forests are being affected by climate change. A Nature study indicates that by 2030, the African jungle will absorb 14 percent less carbon dioxide than it did from around 2005–2010, and will absorb none at all by 2035.
Eastern Africa
Situated almost entirely in the tropics, rainfall in Eastern Africa is dominated by the seasonal migration of the tropical-rain band. Eastern Africa is characterized by high spatio-temporal rainfall variability as it spans over 30 degrees of latitude (across the equator). It has influences from both the Indian and Atlantic Oceans, and has major geographic features (highlands) as well as inland water bodies such as Lake Victoria. Therefore the rainfall seasonality varies from a single wet season per year in July–August in parts of the northwest (including Ethiopia and South Sudan, which are meteorologically more connected to West Africa, with the West African monsoon bringing the rains) to a single wet season per year in December – February in the south (over Tanzania), with many areas close to the equator having two rainy seasons per year, approximately in March–May (the "Long Rains") and October to December (the "Short Rains"). Fine-scale variability in rainfall seasonality is often linked to orography and lakes. Inter-annual variability can be large and known controls include variations in Sea surface temperatures (SSTs) of different ocean basins, large-scale atmospheric modes of variability such as the Madden–Julian Osciliation (MJO) and tropical cyclones. The Long Rains are the main crop-growing season in the region. Interannual predictability of this season is low compared to the Short Rains, and recent drying contrasts with climate projections of a wetter future (the "East African climate paradox".).
Eastern Africa has witnessed frequent and severe droughts in recent decades, as well as devastating floods. Trends in rainfall since the 1980s show a general decrease in March – May (MAM) seasonal rains with a slight increase during June – September (JJAS) and October – December (OND) rains, although there appears to have been a recent recovery in the MAM rains. In the future, both rainfall and temperature are projected to change over Eastern Africa. Recent studies on climate projections suggest that average temperature might increase by about 2–3 °C by the middle of the century and 2–5 °C at the end of the century. This will depend on emission scenarios as well as on how the real climate responds compared with the range of possible outcomes shown by models. Climate model projections tend to show an increase in rainfall, particularly during OND season, which is also projected to occur later. This delay in the short rain season, has been linked to the deepening of the Saharan Heat Low under climate change. It should be noted, however that some models predict decreasing rainfall, and for some regions and seasons the very largest rainfall increases predicted have been shown to involve implausible mechanisms due to systematic model errors. In addition, changes of aerosols provide a forcing of rainfall change that is not captured in many assessments of climate projections.
The contrast of the drying trend of MAM (long rains) rainfall in equatorial Eastern Africa, with most models predicting a wetting in the future has been labelled the "East African climate change paradox", although there has been some recent recovery in the rainfall. Studies have shown that the drying trend is unlikely to be purely natural, but may be driven by factors such as aerosols rather than greenhouse gases, further research is needed. The drying has been shown to have been caused by a shorter rainy season, and linked to deepening of the Arabian Heat Low.
Consistent with the uncertainty in rainfall projections, changes in rainy seasons onset are uncertain in equatorial Eastern Africa, although many models predict a later and wetter short rains. The Indian Ocean Dipole (IOD) is known to provide a strong control on inter-annual variability in the short rains, and studies show that extreme IODs may increase under climate change.
Globally, climate change is expected to lead to intensification of rainfall, as extreme rainfall increases at a faster rate with warming than total rainfall does. Recent work shows that across Africa global models are expected to under-estimate the rate of change of this rainfall intensification, and changes in rainfall extremes may be much more widespread than those predicted by global models.
Southern parts of Eastern Africa receive most of their rainfall in a single rainy season during the southern hemisphere's winter: over Tanzania seasonal rainfall is projected to increase under future climate change, although there is uncertainty. Further south, over Mozambique, a shorter season due to a later onset is projected under future climate change, again with some uncertainty.
North Africa
West Africa and the Sahel
The West African region can be divided into four climatic sub-regions namely the Guinea Coast, Soudano-Sahel, Sahel (extending eastward to the Ethiopian border) and the Sahara, each with different climatic conditions. The seasonal cycle of rainfall is mainly driven by the south-north movement of the Inter-Tropical Convergence Zone (ITCZ) which is characterised by the confluence between moist southwesterly monsoon winds and the dry northeasterly Harmattan.
Based on the inter-annual rainfall variability, three main climatic periods have been observed over the Sahel: the wet period from 1950 to the early 1960s followed by a dry period from 1972 to 1990 and then the period from 1991 onwards which has seen a partial rainfall recovery. During the dry period, the Sahel experienced a number of particularly severe drought events, with devastating effects. The recent decades, have also witnessed a moderate increment in annual rainfall since the beginning of 1990s. However, total annual rainfall remains significantly below that observed during the 1950s.
Some have identified the two recent decades as a recovery period. Others refer to this as a period of "hydrological intensification" with much of the annual rainfall increase coming from more severe rain events and sometimes flooding rather than more frequent rainfall, or similarly other works underline the continuity of the drought even though the rainfall has increased. Since 1985, 54 percent of the population has been affected by five or more floods in the 17 Sahel region countries. In 2012, severe drought conditions in the Sahel were reported. Governments in the region responded quickly, launching strategies to address the issue.
The region is projected to experience changes in rainfall regime, with climate models suggesting that decreases in wet season rainfall are more likely in the western Sahel, and increases more likely in the central to east Sahel, although opposite trends cannot yet be ruled out. These trends will affect the frequency and severity of floods, droughts, desertification, sand and dust storms, desert locust plagues and water shortages.
However, irrespective of the changes in seasonal mean rain, the most intense storms are expected to become more intense, amplifying flood frequency. Enhanced carbon emissions and global warming may also lead to an increase in dry spells especially across the Guinea Coast associated with a reduction of the wet spells under both 1.5 °C and 2 °C global warming level.
Fifteen percent of Sahel region population has also experienced a temperature increase of more than 1 °C from 1970 to 2010. The Sahel region, in particular, will experience higher average temperatures over the course of the 21st century and changes in rainfall patterns, according to the Intergovernmental Panel on Climate Change (IPCC).
Southern Africa
Adaptation
To reduce the impacts of climate change on African countries, adaptation measures are required at multiple scales – ranging from local to national and regional levels. The first generation of adaptation projects in Africa can be largely characterized as small-scale in nature, focused on targeted investments in agriculture and diffusion of technologies to support adaptive decision-making. More recently, programming efforts have re-oriented towards larger and more coordinated efforts, tackling issues that spanning multiple sectors. According to a 2023 study, 59% of African banks have a climate change policy in place, with another 22% planning to implement one. 65% of banks presently consider climate risk when evaluating new clients or projects, with another 23% expecting to do so in the future.
Improved weather forecasting technology in sub-Saharan Africa is important to inform the response to climate change, to aid decision-making associated with adaptation to climate change for example.
During the 21st Conference of the Parties (COP) in 2015, African heads of state launched the Africa Adaptation Initiative (AAI). The AAI's steering committee is composed of the African Ministerial Conference on Environment (AMCEN) Bureau and the chair of the African Group of Negotiators (AGN).
The Africa Adaptation Initiative is also supported by the European Union. The European Union has partnered with the African Union on the promotion of sustainable resources management, environmental resilience, and climate change mitigation
At the regional level, regional policies and actions in support of adaptation across Africa are still in their infancy. The IPCC's Fifth Assessment Report (AR5) highlights examples of various regional climate change action plans, including those developed by the Southern African Development Community (SADC) and Lake Victoria Basin Committee. At the national level, many early adaptation initiatives were coordinated through National Adaptation Programmes of Action (NAPAs) or National Climate Change Response Strategies (NCCRS). Implementation has been slow however, with mixed success in delivery. Integration of climate change with wider economic and development planning remains limited but growing.
At the subnational level, many provincial and municipal authorities are also developing their own strategies, for example the Western Cape Climate Change Response Strategy. Yet, levels of technical capacity and resources available to implement plans are generally low. There has been considerable attention across Africa given to implementing community-based adaptation projects. There is broad agreement that support to local-level adaptation is best achieved by starting with existing local adaptive capacity, and engaging with indigenous knowledge and practices.
The IPCC highlights a number of successful approaches to promote effective adaptation in Africa, outlining five common principles. These include:
Enhancing support for autonomous forms of adaptation;
Increasing attention to the cultural, ethical, and rights considerations of adaptation (especially through active participation of women, youth, and poor and vulnerable people in adaptation activities);
Combining "soft path" options and flexible and iterative learning approaches with technological and infrastructural approaches (including integration of scientific, local, and indigenous knowledge in developing adaptation strategies)
Focusing on enhancing resilience and implementing low-regrets adaptation options; and
Building adaptive management and encouraging process of social and institutional learning into adaptation activities.
The World Health Organization's report "Adaptation to Climate Change in Africa Plan of Action for the Health Sector 2012–2016" is intended to "provide a comprehensive and evidence-based coordinated response of the health sector to climate change adaptation needs of African countries in order to support the commitments and priorities of African governments." The action plan includes goals like scaling up public health activities, coordinating efforts on an international scale, strengthening partnerships and collaborative efforts, and promoting research on both the effects of climate change as well as effective measures taken in local communities to mitigate climate change consequences.
According to the International Monetary Fund (IMF), Sub-Saharan Africa requires $30–$50 billion in additional financing each year to adapt to the effects of climate change.
Climate financing in the Middle East and North Africa totaled $32.6 billion (2% of the world total) in 2019/2020, while climate investment in Sub-Saharan Africa was $43.8 billion (3% of the global total).
According to the European Investment Bank's Banking in Africa study 2021, African institutions are becoming more conscious of the need to address the dangers posed by climate change and are beginning to capitalize on possibilities in green financing. For example, 54% of questioned banks in the study saw climate change as a strategic concern, and more than 40% had people focusing on climate-related fronts. Sub-Saharan African banks are growing their digital offerings, which has been expedited by the COVID-19 pandemic. The majority of the banks surveyed said that the pandemic has accelerated the speed of digital transformation, and that this shift will be permanent.
The poor and vulnerable are most susceptible, with migrant workers, refugees, and other marginalised groups likely to suffer the most. GDP per capita is not likely to rebound to 2019 levels until 2024, with risks tilting to the downside, and the crisis has reversed a predicted drop in the number of poor people, according to the IMF.
In comparison to pre-crisis forecasts, this might result in an additional 30 million people in Sub-Saharan Africa living in extreme poverty by 2021, as well as an additional nine million in the Middle East and North Africa (MENA) area.
As of 2023, about a third of all African climate funding flows to five major markets: Morocco (7% of African climate investment in 2019/2022), Nigeria (7%), Kenya (7%), Ethiopia (6%), and South Africa (5%).
Over the last decade, worldwide greenfield foreign direct investment has declined at a 3% annual rate, with Africa's global contribution dropping from 12% in 2017 to less than 6% in 2021.
Northern Africa adaptation measures
Key adaptations in northern Africa relate to increased risk of water scarcity (resulting from a combination of climate change affecting water availability and increasing demand). Reduced water availability, in turn, interacts with increasing temperatures to create need for adaptation among rainfed wheat production and changing disease risk (for example from leishmaniasis). Most government actions for adaptation centre on water supply side, for example through desalination, inter-basin transfers and dam construction. Migration has also been observed to act as an adaptation for individuals and households in northern Africa. Like many regions, however, examples of adaptation action (as opposed to intentions to act, or vulnerability assessments) from north Africa are limited – a systematic review published in 2011 showed that only 1 out of 87 examples of reported adaptations came from North Africa.
Western Africa adaptation measures
Water availability is a particular risk in Western Africa, with extreme events such as drought leading to humanitarian crises associated with periodic famines, food insecurity, population displacement, migration and conflict and insecurity. Adaptation strategies can be environmental, cultural/agronomic and economic.
Adaptation strategies are evident in the agriculture sector, some of which are developed or promoted by formal research or experimental stations. Indigenous agricultural adaptations observed in northern Ghana are crop-related, soil-related or involve cultural practices. Livestock-based agricultural adaptations include indigenous strategies such as adjusting quantities of feed to feed livestock, storing enough feed during the abundant period to be fed to livestock during the lean season, treating wounds with solution of certain barks of trees, and keeping local breeds which are already adapted to the climate of northern Ghana; and livestock production technologies to include breeding, health, feed/nutrition and housing.
The choice and adoption of adaptation strategies is variously contingent on demographic factors such as the household size, age, gender and education of the household head; economic factors such as income source; farm size; knowledge of adaptation options; and expectation of future prospects.
Eastern Africa adaptation measures
In Eastern Africa adaptation options are varied, including improving use of climate information, actions in the agriculture and livestock sector, and in the water sector.
Making better use of climate and weather data, weather forecasts, and other management tools enables timely information and preparedness of people in the sectors such as agriculture that depend on weather outcomes. This means mastering hydro-meteorological information and early warning systems. It has been argued that the indigenous communities possess knowledge on historical climate changes through environmental signs (e.g. appearance and migration of certain birds, butterflies etc.), and thus promoting of indigenous knowledge has been considered an important adaptation strategy.
Adaptation in the agricultural sector includes increased use of manure and crop-specific fertilizer, use of resistant varieties of crops and early maturing crops. Manure, and especially animal manure is thought to retain water and have essential microbes that breakdown nutrients making them available to plants, as compared to synthetic fertilizers that have compounds which when released to the environment due to over-use release greenhouse gases. One major vulnerability of the agriculture sector in Eastern Africa is the dependence on rain-fed agriculture. An adaptation solution is efficient irrigation mechanisms and efficient water storage and use. Drip irrigation has especially been identified as a water-efficient option as it directs the water to the root of the plant with minimal wastage. Countries like Rwanda and Kenya have prioritized developing irrigated areas by gravity water systems from perennial streams and rivers in zones vulnerable to prolonged droughts. During heavy rains, many areas experience flooding resulting from bare grounds due to deforestation and little land cover. Adaptation strategies proposed for this is promoting conservation efforts on land protection, by planting indigenous trees, protecting water catchment areas and managing grazing lands through zoning.
For the livestock sector, adaptation options include managing production through sustainable land and pasture management in the ecosystems. This includes promoting hay and fodder production methods e.g. through irrigation and use of waste treated water, and focusing on investing in hay storage for use during dry seasons. Keeping livestock is considered a livelihood rather than an economic activity. Throughout Eastern Africa Countries especially in the ASALs regions, it is argued that promoting commercialisation of livestock is an adaptation option. This involves adopting economic models in livestock feed production, animal traceability, promoting demand for livestock products such as meat, milk and leather and linking to niche markets to enhance businesses and provide disposable income.
In the water sector, options include efficient use of water for households, animals and industrial consumption and protection of water sources. Campaigns such as planting indigenous trees in water catchment areas, controlling human activities near catchment areas especially farming and settlement have been carried out to help protect water resources and avail access to water for communities especially during climatic shocks.
Comoros – "NAPA is the operational extension of the Poverty Reduction Strategy Paper (PRSP), as it includes among its adaptation priorities, agriculture, fishing, water, housing, health, but also tourism, in an indirect way, through the reconstitution of basin slopes and the fight against soils erosion, and therefore the protection of reefs by limiting the silting up by terrigenous contributions."
Kenya gazetted the Climate Change Act, 2016 which establishes an authority to oversee development, management, implementation and regulation of mechanisms to enhance climate change resilience and low carbon development for sustainable development, by the National and County Governments, the private sector, civil society, and other actors. Kenya has also developed the National Climate Change Action Plan (NCCAP 2018–2022 ) which aims to further the country's development goals by providing mechanisms and measures to achieve low carbon climate-resilient development in a manner that prioritizes adaptation.
Central Africa adaptation measures
Angola – "The objective of the National Adaptation Programs of Action are to identify and communicate the urgent and immediate needs of the country regarding climate change adaptation, to increase Angola's resilience to climate variabilities and to climate change to ensure achievement of Poverty reduction programs, sustainable development objectives and the Millennium Development Goals pursued by the Government."
Southern Africa adaptation measures
There have been several initiatives at local (site-specific), local, national and regional scales aimed at strengthening to climate change. Some of these are: The Regional Climate Change Programme (RCCP), SASSCAL, ASSAR, UNDP Climate Change Adaptation, RESILIM, FRACTAL. South Africa implemented the Long-Term adaptation Scenarios Flagship Research Programme (LTAS) from April 2012 to June 2014. This research also produced factsheets and a technical report covering the SADC region entitled "Climate Change Adaptation: Perspectives for the Southern African Development Community (SADC)".
Madagascar – the priority sectors for adaptation are: agriculture and livestock, forestry, public health, water resources and coastal zones.
Malawi – The NAPA identifies the following as high priority activities for adaptation: "Improving community resilience to climate change through the development of sustainable rural livelihoods, Restoring forests in the Upper and Lower Shire Valleys catchments to reduce siltation and associated water flow problems, Improving agricultural production under erratic rains and changing climatic conditions, Improving Malawi's preparedness to cope with droughts and floods, and Improving climate monitoring to enhance Malawi's early warning capability and decision making and sustainable utilisation of Lake Malawi and lakeshore areas resources". And according to the World Bank's Country Climate and Development Report (CCDR) for Malawi, can "take steps to jumpstart investments in climate-resilient infrastructure and halt land degradation and forest loss to improve agriculture productivity and carbon capture" "
Mauritius – adaptation should address the following priority areas: coastal resources, agriculture, water resources, fisheries, health and well-being, land use change and forestry and biodiversity.
Mozambique – "The proposed adaptation initiatives target various areas of economic and social development, and outline projects related to the reduction of impacts to natural disasters, the creation of adaptation measures to climate change, fight against soil erosion in areas of high desertification and coastal zones, reforestation and the management of water resources.""
Rwanda has developed the National Adaptation Programme of Action (NAPA 2006) which contains information to guide national policy-makers and planners on priority vulnerabilities and adaptations in important economic sectors. The country has also developed sector based policies on adaptation to climate change such as the Vision 2020, the National Environmental Policy and the Agricultural Policy among others.
Tanzania – Tanzania has outlined priority adaptation measures in their NAPA, and various national sector strategies and research outputs. The NAPA has been successful at encouraging climate change mainstreaming into sector policies in Tanzania; however, the cross-sectoral collaboration crucial to implementing adaptation strategies remains limited due to institutional challenges such as power imbalances, budget constraints and an ingrained sectoral approach. Most of the projects in Tanzania concern agriculture and water resource management (irrigation, water saving, rainwater collection); however, energy and tourism also play an important role.
Zambia – "The NAPA identifies 39 urgent adaptation needs and 10 priority areas within the sectors of agriculture and food security (livestock, fisheries and crops), energy and water, human health, natural resources and wildlife."
Zimbabwe – "The other strategic interventions by the NAP process will be: Strengthening the role of private sector in adaptation planning, Enhancing of the capacity of Government to develop bankable projects through trainings, Improving management of background climate information to inform climate change planning, Crafting a proactive resource-mobilization strategy for identifying and applying for international climate finance as requests for funds are primarily reactive at present, focusing on emergency relief rather than climate change risk reduction, preparedness and adaptation, Developing a coordinated monitoring and evaluation policy for programs and projects, as many institutions within the government do not currently have a systematic approach to monitoring and evaluation."
Lesotho – "The key objectives of the NAPA process entail: identification of communities and livelihoods most vulnerable to climate change, generating a list of activities that would form a core of the national adaptation program of action, and to communicate the country's immediate and urgent needs and priorities for building capacity for adaptation to climate change.""
Namibia – the critical themes for adaptation are "Food security and sustainable biological resource base, Sustainable water resources base, Human health and well being and Infrastructure development.
South Africa has adopted in August 2020 its National Climate Change Adaptation Strategy, which "acts as a common reference point for climate change adaptation efforts in South Africa, and it provides a platform upon which national climate change adaptation objectives for the country can be articulated so as to provide overarching guidance to all sectors of the economy"
Society and culture
Inequality in climate research
Even though Africa is going to be one of the most affected continents from climate change, systematic inequity and other biases related to scientific research and funding mean that very little of the published science about climate change and climate research funding is for African scientist. An analysis of research money from 1990 to 2020 for climate change, found that 78% of research money for research on climate change in Africa was spent in European and North American institutions and more was spent for former British colonies than other countries. This pattern of parachute science, in turn both prevents local researchers from doing groundbreaking work, because they do not have the funding for experimental activities and reduces investment in local researchers ideas and in topics important to the Global South, such as climate change adaptation.
Accurate sustainability evaluations are challenging due to a lack of sustainable investment frameworks, as well as data and managerial capability restrictions. Currently, fewer than half of Africa's top pension funds report information on sustainability policies and execution.
The United Nations Conference on Trade and Development - International Standards of Accounting and Reporting (UNCTAD-ISAR) founded the African Regional Partnership for Sustainability and SDG Reporting in 2022. The collaboration has 53 members as of March 2023, including national corporate social responsibility networks and/or ministries from 27 African nations.
| Physical sciences | Climate change | Earth science |
5067669 | https://en.wikipedia.org/wiki/Potassium%20sulfide | Potassium sulfide | Potassium sulfide is an inorganic compound with the formula K2S. The colourless solid is rarely encountered, because it reacts readily with water, a reaction that affords potassium hydrosulfide (KSH) and potassium hydroxide (KOH). Most commonly, the term potassium sulfide refers loosely to this mixture, not the anhydrous solid.
Structure
It adopts "antifluorite structure," which means that the small K+ ions occupy the tetrahedral (F−) sites in fluorite, and the larger S2− centers occupy the eight-coordinate sites. Li2S, Na2S, and Rb2S crystallize similarly.
Synthesis and reactions
It can be produced by heating K2SO4 with carbon (coke):
K2SO4 + 4 C → K2S + 4 CO
In the laboratory, pure K2S may be prepared by the reaction of potassium and sulfur in anhydrous ammonia.
Sulfide is highly basic, consequently K2S completely and irreversibly hydrolyzes in water according to the following equation:
K2S + H2O → KOH + KSH
For many purposes, this reaction is inconsequential since the mixture of SH− and OH− behaves as a source of S2−. Other alkali metal sulfides behave similarly.
Use in fireworks
Potassium sulfides are formed when black powder is burned and are important intermediates in many pyrotechnic effects, such as senko hanabi and some glitter formulations.
| Physical sciences | Sulfide salts | Chemistry |
5068752 | https://en.wikipedia.org/wiki/Trigonotarbida | Trigonotarbida | The order Trigonotarbida is a group of extinct arachnids whose fossil record extends from the late Silurian to the early Permian (Pridoli to Sakmarian). These animals are known from several localities in Europe and North America, as well as a single record from Argentina. Trigonotarbids can be envisaged as spider-like arachnids, but without silk-producing spinnerets. They ranged in size from a few millimetres to a few centimetres in body length and had segmented abdomens (opisthosoma), with the dorsal exoskeleton (tergites) across the backs of the animals' abdomens, which were characteristically divided into three or five separate plates. Probably living as predators on other arthropods, some later trigonotarbid species were quite heavily armoured and protected themselves with spines and tubercles. About seventy species are currently known, with most fossils originating from the Carboniferous coal measures.
Historical background
The first trigonotarbid was described in 1837 from the coal measures of Coalbrookdale in England by the famous English geologist Dean William Buckland. He believed it to be a fossil beetle and named it Curculoides prestvicii. A much better preserved example was later discovered from Coseley near Dudley; also in the English West Midlands conurbation. Described in 1871 by Henry Woodward, he correctly identified it as an arachnid and renamed it Eophrynus prestvicii—whereby the genus name comes from (, meaning 'dawn'), and Phrynus, a genus of living whip spider (Amblypygi). Woodward subsequently described another trigonotarbid, Brachypyge carbonis, from the coal measures of Mons in Belgium; although this fossil is known only from its abdomen and was initially mistaken for those of a crab.
A new arachnid order
In 1882, the German zoologist Ferdinand Karsch described a number of fossil arachnids from the coal measures of Neurode in Silesia (now Poland), including one he named Anthracomartus voelkelianus in honour of Herr Völkel, the foreman of the mine where it was discovered. This species was raised to a new, extinct, arachnid order which Karsch called Anthracomarti. The name is derived from (), the Greek word for coal. A number of other fossils which would eventually be placed in Trigonotarbida were discovered around this time. Hanns Bruno Geinitz described Kreischeria wiedei from the coal measures of Zwickau in Germany, although he interpreted it as a fossil pseudoscorpion. Johann Kušta described Anthracomartus krejcii from Rakovník in the Czech Republic, and published further descriptions in a number of subsequent papers. In 1884, Samuel Hubbard Scudder described Anthracomartus trilobitus from Fayetteville, Arkansas—the first trigonotarbid from North America.
Relationships
Early studies tended to confuse trigonotarbids with other living or extinct groups of arachnids; particularly harvestmen (Opiliones). Petrunkevitch's division of the trigonotarbids into two, unrelated, orders was noted above. In detail, he divided the arachnids into suborders based on the width of the division between the two parts of the body (the prosoma and opisthosoma). Anthracomartida and another extinct order, Haptopoda, were grouped into a subclass Stethostomata defined by a broad division of the body and downward-hanging mouthparts. Trigonotarbida was placed in its own subclass Soluta and defined as having a division of the body which was variable in width. Petrunkevitch's scheme was largely followed in subsequent studies of fossil arachnids.
Pantetrapulmonata
In the 1980s, Bill Shear and colleagues carried out an important study on well preserved Mid Devonian trigonotarbids from Gilboa, New York. They questioned whether it was appropriate to define a group of animals on a variable character state and carried out the first cladistic analysis of fossil and living arachnids. They showed that trigonotarbids are closely related to a group of arachnids which have gone under various names (Caulogastra, Arachnidea, etc.), but for which the name Tetrapulmonata has become most widespread. Members of the Tetrapulmonata include spiders (Araneae), whip spiders (Amblypygi), whip scorpions (Uropygi) and shorttailed whipscorpion (Schizomida) and, together with trigonotarbids, share characters like two pairs of book lungs and similar mouthparts with fangs operating rather like a pocket knife. In a 2007 study of arachnid relationships, the Shear et al. hypothesis was largely supported and a group Pantetrapulmonata was proposed which comprises Trigonotarbida + Tetrapulmonata. This has since been corroborated in more recent cladistic analyses.
Trigonotarbids and ricinuleids
In 1892, Ferdinand Karsch suggested that the rare and rather bizarre-looking ricinuleids (Ricinulei) were the last living descendants of the trigonotarbids. A similar hypothesis was reintroduced by Dunlop, who pointed out distinct similarities and possible sister group relationship between these arachnid groups. Both have opisthosomal tergites divided into median and lateral plates and both have a complicated coupling mechanism between the prosoma and the opisthosoma which 'locks' the two halves of the body together. Although cladistic analysis has tended to recover ricinuleids in their traditional position closely related to mites and ticks, further discoveries have revealed that the tip of the pedipalp ends in a small claw in both trigonotarbids and ricinuleids. If the hypothesis is true, ricinuleids, despite the lack of tetrapulmonate key characters (e.g. book lungs), may represent part of the pantetrapulmonate clade alongside trigonotarbids as well.
Internal relationships
The first cladistic analysis of the trigonotarbids was published in 2014. This recovered the families Anthracomartidae, Anthracosironidae, and Eophrynidae as monophyletic. In contrast Trigonotarbidae, Aphantomartidae, Palaeocharinidae, and Kreischeriidae were not. Two clades were consistently recovered with strong support—(Palaeocharinus (Archaeomartidae + Anthracomartidae)), and Lissomartus as sister group the 'eophrynid assemblage' (Aphantomartus (Alkenia (Pseudokreischeria (Kreischeria (Eophrynus + Pleophrynus))))).
Description
Trigonotarbids superficially resemble spiders, but can be easily recognised by having tergites on the dorsal side of the opisthosoma divided into median and lateral plates. This character is shared with ricinuleids (Ricinulei) (see also Ricinulei#Relationships). As in other arachnids, the body is divided into a prosoma (or cephalothorax) and opisthosoma (or abdomen). Body length ranges from a couple of millimetres up to about .
Prosoma
The prosoma is covered by the carapace and always bears a pair of median eyes. In the probably basal families Palaeocharinidae, Anthracomartidae—and perhaps also Anthracosironidae—there is an additional pair of lateral eye tubercles which, at least in palaeocharinids, appear to have borne a series of individual lenses. In this sense palaeocharinids seem to be in the process of reducing a compound eye. Anterior margin of the carapace protrude into a projection referred to as clypeus.
The chelicerae are of the "pocket-knife" type consisting of a basal segment and a sharp, curving fang. The chelicerae are described as paleognathic: the fangs are held parallel to one another, like those of mesothele and mygalomorph spiders, but the chelicerae hang downwards like those of araneomorph spiders. There is no evidence in well-preserved fossils for the opening of a venom gland, thus trigonotarbids were probably not venomous. The chelicerae may have been slightly retractable into the prosoma. Well-preserved palaeocharinids show evidence for a small, slit-like mouth with an upper lip (a labrum or rostrum) and a lower lip (or labium). Inside the mouth there is some sort of filtering system formed from hairs or platelets which strongly suggests that trigonotarbids (like spiders and many other arachnids) could eat only preorally digested, liquified prey.
The pedipalps have the typical arachnid structure with a coxa, trochanter, femur, patella, tibia and tarsus. They are pediform, i.e. they look like small legs and were not highly modified. There is no evidence for a special sperm transfer device as in the modified palpal organ of male spiders. In at least the palaeocharinids and anthracomartids the tip of the pedipalp is modified into a small chela (claw) formed from the tarsal claw (or apotele) and a projection from the tarsus. As mentioned above, a very similar arrangement is seen at the end of the pedipalp in Ricinulei.
The walking legs again follow the typical arachnid plan with a coxa, trochanter, femur, patella, tibia, metatarsus and tarsus. The coxae surround a single sternum. In well preserved palaeocharinids there is a ring, or annulus, around the trochanter–femur joint which may be the remains of an earlier leg segment. The legs are largely unmodified, although in Anthracosironidae the forelegs are quite large and spiny, presumably to help catch prey. The legs end in three claws, two large ones and a smaller median claw.
Opisthosoma
The opisthosoma is largely suboval in outline with a flatten dorsal surface. It compose of 12 segments, with some of them had undergone degrees of fusion or reduction, hence the previous misinterpretation of around 8 to 11 segments. Tergite of the first segment partially covered by the posterior margin of preceding carapace, forming a complicated coupling mechanism known as 'locking ridge'. Tergites of segment 2 to 8 (segment 9 in some species) were all laterally divided into 3 (one median and two lateral) plates, with those of segment 2 and 3 fused to each other in most species. However, the corresponding tergites of the family Anthracomartidae are further subdivided into 5 plates. The last 3 segments are usually only visible from the ventral side, with the 2 final segments constricted into a tiny ring-like section known as pygidium.
Ventral side of opisthosomal segment 2 to 9 covered by series of lung-bearing opercula (2 and 3) and curved sternites (4 to 9). The first segment apparently lacking any ventral plates. Just like other lung-bearing arachnids (scorpion and tetrapulmonate), the book lungs of trigonotarbids formed by layers of trabecula-bearing lamellae, which is a feature adapted to a terrestrial, air-breathing lifestyle. A pair of ventral sacs located between the posterior operculum and following sternite had been observed in some species.
Paleobiology
In July 2014 scientists used computer-based techniques to re-create a possible walking gait for the animal. A subsequent review article suggested by comparison with mites, with presumably similar lifestyle and environment, a metachronal rather than alternating leg coordination was more likely. Subsequent work by the researchers behind the initial publication used simulation approaches to assess the efficiency of a range of gaits using an updated trigonotarbid model.
Included taxa
As of 2020, 70 valid species had been included under Trigonotarbida as follows:
plesion taxa
Palaeotarbus Dunlop, 1999
Palaeotarbus jerami (Dunlop, 1996) – Late Silurian, England
Palaeocharinidae Hirst, 1923
Aculeatarbus Shear, Selden & Rolfe, 1987
Aculeatarbus depressus Shear, Selden & Rolfe, 1987 – Mid Devonian, United States
Gelasinotarbus Shear, Selden & Rolfe, 1987
Gelasinotarbus bifidus Shear, Selden & Rolfe, 1987 – Mid Devonian, United States
Gelasinotarbus bonamoae Shear, Selden & Rolfe, 1987 – Mid Devonian, United States
Gelasinotarbus heptops Shear, Selden & Rolfe, 1987 – Mid Devonian, United States
Gelasinotarbus reticulatus Shear, Selden & Rolfe, 1987 – Mid Devonian, United States
Gigantocharinus Shear, 2000
Gigantocharinus szatmaryi Shear, 2000 – Late Devonian, United States
Gilboarachne Shear, Selden & Rolfe, 1987
Gilboarachne griersoni Shear, Selden & Rolfe, 1987 – Mid Devonian, United States
Palaeocharinus Hirst, 1923
Palaeocharinus calmani Hirst, 1923 – Early Devonian, Scotland
Palaeocharinus hornei Hirst, 1923 – Early Devonian, Scotland
Palaeocharinus kidstoni Hirst, 1923 – Early Devonian, Scotland
Palaeocharinus rhyniensis Hirst, 1923 – Early Devonian, Scotland
Palaeocharinus scourfieldi Hirst, 1923 – Early Devonian, Scotland
Palaeocharinus tuberculatus Fayers, Dunlop & Trewin, 2005 – Early Devonian, Scotland
Spinocharinus Poschmann & Dunlop, 2011
Spinocharinus steinmeyeri Poschman & Dunlop, 2011 - Devonian, Bürdenbach
Archaeomartidae Haase, 1890
Archaeomartus Størmer, 1970
Archaeomartus levis Størmer, 1970 - Devonian, Alken an der Mosel
Anthracomartidae Haase, 1890
synonyms
= Promygalidae Frič, 1904
= Brachypygidae Pocock, 1911
= Coryphomartidae Petrunkevitch, 1945
= Pleomartidae Petrunkevitch, 1945
Anthracomartus Karsch, 1882
synonyms
= Brachylycosa Frič, 1904
= Cleptomartus Petrunkevitch, 1949
= Coryphomartus Petrunkevitch, 1945
= Cryptomartus Petrunkevitch, 1945
= Oomartus Petrunkevitch, 1953
= Perneria Frič, 1904
= Pleomartus Petrunkevitch, 1945
= Promygale Frič, 1901
Anthracomartus bohemica (Frič, 1901) – Late Carboniferous, Czech Republic
Anthracomartus carcinoides (Frič, 1901) – Late Carboniferous, Czech Republic
synonyms
= Promygale rotunda Frič, 1901
= Perneria salticoides Frič, 1904
Anthracomartus elegans Frič, 1901 – Late Carboniferous, Czech Republic
Anthracomartus hindi Pocock, 1911 – Late Carboniferous, England
synonyms
= Cleptomartus hangardi Guthörl, 1965
= Cryptomartus meyeri Guthörl, 1964
= Cleptomartus planus Petrunkevitch, 1949
= Cryptomartus rebskei Brauckmann, 1984
Anthracomartus granulatus Frič, 1904 – Late Carboniferous, Poland
Anthracomartus janae (Opluštil, 1986) – Late Carboniferous, Czech Republic
Anthracomartus kustae Petrunkevitch, 1953 – Late Carboniferous, Czech Republic
Anthracomartus minor Kušta, 1884 – Late Carboniferous, Czech Republic
synonym
= Anthracomartus socius Kušta, 1888
Anthracomartus nyranensis (Petrunkevitch, 1953) – Late Carboniferous, Czech Republic
Anthracomartus palatinus Ammon, 1901 – Late Carboniferous, Germany
Anthracomartus priesti Pocock, 1911 – Late Carboniferous, England
synonyms
= Anthracomartus denuiti Pruvost, 1922
= Cleptomartus plautus Petrunkevitch, 1949
Anthracomartus radvanicensis (Opluštil, 1985) – Late Carboniferous, Czech Republic
Anthracomartus triangularis Petrunkevitch, 1913 – Late Carboniferous, Canada
Anthracomartus trilobitus Scudder, 1884 – Late Carboniferous, United States
Anthracomartus voelkelianus Karsch, 1882 – Late Carboniferous, Poland
Brachypyge Woodward, 1878
Brachypyge carbonis Woodward, 1878 – Late Carboniferous, Belgium
Maiocercus Pocock, 1911
Maiocercus celticus (Pocock, 1902) – Late Carboniferous, Europe
synonym
= Maiocercus orbicularis Gill, 1911
Anthracosironidae Pocock, 1903
Anthracosiro Pocock, 1903
Anthracosiro fritschii Pocock, 1903 – Late Carboniferous, Europe
synonym
= Anthracosiro elongatus Waterlot, 1934
Anthracosiro woodwardi Pocock, 1903 – Late Carboniferous, Europe
synonyms
= Anthracosiro corsini Pruvost, 1926
= Anthracosiro latipes Gill, 1909
Arianrhoda Dunlop & Selden, 2004
Arianrhoda bennetti Dunlop & Selden, 2004 – Early Devonian, Wales
Vratislavia Frič, 1904
Vratislavia silesica (Roemer, 1878) - Carboniferous, Silesia
Trigonotarbidae Petrunkevitch, 1949
Trigonotarbus Pocock, 1911
Trigonotarbus arnoldi Petrunkevitch, 1955 – Late Carboniferous, France
Trigonotarbus johnsoni Pocock, 1911 – Late Carboniferous, England
Trigonotarbus stoermeri Schultka, 1991 – Early Devonian, Germany
Lissomartidae Dunlop, 1995
Lissomartus Petrunkevitch, 1949
Lissomartus carbonarius (Petrunkevitch, 1913) – Late Carboniferous, United States
Lissomartus schucherti (Petrunkevitch, 1913) – Late Carboniferous, United States
Aphantomartidae Petrunkevitch, 1945
synonym
= Trigonomartidae Petrunkevitch, 1949
Alkenia Størmer, 1970
Alkenia mirabilis Størmer, 1970 - Devonian, Alken an der Mosel
Aphantomartus Pocock, 1911
synonyms
= Trigonomartus Petrunkevitch, 1913
= Phrynomartus Petrunkevitch, 1945a
Aphantomartus areolatus Pocock, 1911 – Early/Late Carboniferous, Europe
synonyms
= Aphantomartus pococki Pruvost, 1912
= Trigonomartus dorlodoti Pruvost, 1930
= Eophrynus waechteri Guthörl, 1938
= ?Trigonomartus pruvosti van der Heide, 1951
= ?Brachylycosa manebachensis Müller, 1957
Aphantomartus ilfeldicus (Scharf, 1924) – Permian, Germany
Aphantomartus pustulatus (Scudder, 1884) – Late Carboniferous, Europe, North America
synonyms
= ?Kreischeria villeti Pruvost, 1912
= Cleptomartus plötzensis Simon, 1971
Kreischeriidae Haase, 1890
Anzinia Petrunkevitch, 1953
Anzinia thevenini (Pruvost, 1919) – Late Carboniferous, France
Gondwanarache Pinto & Hünicken, 1980
Gondwanarache argentinensis Pinto & Hünicken, 1980 – Late Carboniferous, Argentina
Hemikreischeria Frič, 1904
Hemikreischeria geinitzi (Thevenin, 1902) – Late Carboniferous, France
Kreischeria Geinitz, 1882
Kreischeria wiedei Geinitz, 1882 – Late Carboniferous, Germany
Pseudokreischeria Petrunkevitch, 1953
Pseudokreischeria pococki (Gill, 1924) – Late Carboniferous, England
synonym
= Eophrynus varius Petrunkevitch, 1949
Eophrynidae Karsch, 1882
synonym
= Hemiphrynidae Frič, 1904
Eophrynus Woodward, 1871
Eophrynus prestvicii (Buckland, 1837) – Late Carboniferous, England
Eophrynus udus Brauckmann, Koch & Kemper, 1985 – Late Carboniferous, Germany
Nyranytarbus Harvey & Selden, 1995
synonym
Hemiphrynus Frič, 1901
Nyranytarbus hofmanni (Frič, 1901) – Late Carboniferous, Czech Republic
Nyranytarbus longipes (Frič, 1901) – Late Carboniferous, Czech Republic
Petrovicia Frič, 1904
Petrovicia proditoria Frič, 1904 – Late Carboniferous, Czech Republic
Planomartus Petrunkevitch, 1953
Planomartus krejcii (Kušta, 1883) – Late Carboniferous, Czech Republic
synonym
= Anthracomartus affinis Kušta, 1885
Pleophrynus Petrunkevitch, 1945a
Pleophrynus verrucosus (Pocock, 1911) – Late Carboniferous, UK, United States
synonym
= Eophrynus warei Dix & Pringle, 1930
= Pleophrynus ensifer Petrunkevitch, 1945a
= Eophrynus jugatus Ambrose & Romano, 1972
Pocononia Petrunkevitch, 1953
Pocononia whitei (Ewing, 1930) – Early Carboniferous, United States
Somaspidion Jux, 1982
Somaspidion hammapheron Jux, 1982
Stenotrogulus Frič, 1904
synonyms
= Cyclotrogulus Frič, 1904
= Pseudoeophrynus Příbyl, 1958
Stenotrogulus salmii (Stur, 1877) – Late Carboniferous, Czech Republic
synonyms
= Cyclotrogulus sturii Frič, 1904 [non Hasse, 1890]
= Pseudoeophrynus ostraviensis Příbyl, 1958
Family uncertain
Aenigmatarbus Poschmann, Dunlop, Bértoux & Galtier, 2016
Aenigmatarbus rastelli Poschmann, Dunlop, Bértoux & Galtier, 2016 - Carboniferous, Graissessac, France
Namurotarbus Poschmann & Dunlop, 2010
Namurotarbus roessleri (Dunlop & Brauckmann, 2006) - Carboniferous, Hagen-Vorhalle
synonyms
= Archaeomartus roessleri Dunlop & Brauckmann, 2006
Permotarbus Dunlop & Rößler, 2013
Permotarbus schuberti Dunlop & Rößler, 2013 Permian, Chemnitz
Tynecotarbus Hradská & Dunlop, 2013
Tynecotarbus tichaveki Hradská & Dunlop, 2013 - Carboniferous, Týnec
incertae sedis
Anthracophrynus Andrée, 1913
Anthracophrynus tuberculatus Andrée, 1913 – Late Carboniferous, Germany
Areomartus Petrunkevitch, 1913
Areomartus ovatus Petrunkevitch, 1913 - Carboniferous, West Virginia
‘Eophrynus’ scharfi Scharf, 1924 – Early Permian, Germany
Aphantomartus Pocock, 1911
Aphantomartus woodruffi (Scudder, 1893) - Carboniferous, Rhode Island
nomina dubia
Anthracomartus buchi (Goldenberg, 1873) – Late Carboniferous, Germany
Anthracomartus hageni (Goldenberg, 1873) – Late Carboniferous, Germany
Elaverimartus pococki Petrunkevitch, 1953 – Late Carboniferous, Scotland
Eurymartus latus Matthew, 1895 – Late Carboniferous, Canada
?Eurymartus spinulosus Matthew, 1895 – Late Carboniferous, Canada
| Biology and health sciences | Prehistoric arachnids | Animals |
5069414 | https://en.wikipedia.org/wiki/Genital%20herpes | Genital herpes | Genital herpes is a herpes infection of the genitals caused by the herpes simplex virus (HSV). Most people either have no or mild symptoms and thus do not know they are infected. When symptoms do occur, they typically include small blisters that break open to form painful ulcers. Flu-like symptoms, such as fever, aching, or swollen lymph nodes, may also occur. Onset is typically around 4 days after exposure with symptoms lasting up to 4 weeks. Once infected further outbreaks may occur but are generally milder.
The disease is typically spread by direct genital contact with the skin surface or secretions of someone who is infected. This may occur during sex, including anal, oral, and manual sex. Sores are not required for transmission to occur. The risk of spread between a couple is about 7.5% over a year. HSV is classified into two types, HSV-1 and HSV-2. While historically HSV-2 was more common, genital HSV-1 has become more common in the developed world. Diagnosis may occur by testing lesions using either PCR or viral culture or blood tests for specific antibodies.
Efforts to prevent infection include not having sex, using condoms, and only having sex with someone who is not infected. Once infected, there is no cure. Antiviral medications may, however, prevent outbreaks or shorten outbreaks if they occur. The long-term use of antivirals may also decrease the risk of further spread.
In 2015, about 846 million people (12% of the world population) had genital herpes. In the United States, more than one in six people between the ages of 14 and 49 have the disease. Women are more commonly infected than men. Rates of disease caused by HSV-2 have decreased in the United States between 1990 and 2010. Complications may rarely include aseptic meningitis, an increased risk of HIV/AIDS if exposed to HIV-positive individuals, and spread to the baby during childbirth resulting in neonatal herpes.
Signs and symptoms
In males, the lesions occur on the glans penis, shaft of the penis or other parts of the genital region, on the inner thigh, buttocks, or anus. In females, lesions appear on or near the pubis, clitoris or other parts of the vulva, buttocks or anus.
Other common symptoms include pain, itching, and burning. Less frequent, yet still common, symptoms include discharge from the penis or vagina, fever, headache, muscle pain (myalgia), swollen and enlarged lymph nodes and malaise. Women often experience additional symptoms that include painful urination (dysuria) and cervicitis. Herpetic proctitis (inflammation of the anus and rectum) is common for individuals participating in anal intercourse.
After 2–3 weeks, existing lesions progress into ulcers and then crust and heal, although lesions on mucosal surfaces may never form crusts. In rare cases, involvement of the sacral region of the spinal cord can cause acute urinary retention and one-sided symptoms and signs of myeloradiculitis (a combination of myelitis and radiculitis): pain, sensory loss, abnormal sensations (paresthesia) and rash. Historically, this has been termed Elsberg syndrome, although this entity is not clearly defined.
Recurrence
After a first episode of herpes genitalis caused by HSV-2, there will be at least one recurrence in approximately 80% of people, while the recurrence rate for herpes genitalis caused by HSV-1 is approximately 50%. Herpes genitalis caused by HSV-2 recurs on average four to six times per year, while that of HSV-1 infection occurs only about once per year.
People with recurrent genital herpes may be treated with suppressive therapy, which consists of daily antiviral treatment using acyclovir, valacyclovir or famciclovir. Suppressive therapy may be useful in those who have at least four recurrences per year but the quality of the evidence is poor. People with lower rates of recurrence will probably also have fewer recurrences with suppressive therapy. Suppressive therapy should be discontinued after a maximum of one year to reassess recurrence frequency.
Transmission
Genital herpes can be spread by viral shedding prior to and following the formation of ulcers. The risk of spread between a couple is about 7.5% over a year (for unprotected sex). The likelihood of transferring genital herpes from one person to another is decreased by external condom use by 50%, by internal condom by 50%, and refraining from sex during an active outbreak. The longer a partner has had the infection, the lower the transmission rate. An infected person may further decrease transmission risks by maintaining a daily dose of antiviral medications. Infection by genital herpes occurs in about 1 in every 1,000 sexual acts.
Prevention
Because herpes simplex virus (HSV) infection is common and not routinely screened for in the general population, complete prevention of the transmission of genital herpes is difficult. To reduce the chance of contracting herpes simplex virus, external condoms for the penis may be used during oral sex, vaginal sex, and anal sex. Internal condoms for the vagina may be used during oral sex or vaginal sex. Internal condoms and external condoms should not be used simultaneously. Dental dams may be used during oral sex involving the vagina or anus. Decreasing the number of sexual partners a person has may also decrease the chance of contracting HSV. People who have sexual relations with others may get tested for HSV. In people who have been diagnosed with genital herpes, transmission to others may be prevented through suppressive antiviral drugs. This option is 90% effective in preventing the transmission of HSV and is a commonly used option for sexual and/or romantic partners or those who plan on becoming pregnant. Those who are aware that they have genital herpes should notify their partner(s).
Screening and diagnosis
Genital herpes may be diagnosed through a physical examination by a doctor or through a herpes simplex virus (HSV) test by sampling fluid within a genital blister or blood for HSV antibodies. Herpes simplex virus testing is recommended for those who have symptoms of herpes or who have a sexual partner who has a herpes infection. There is currently no recommendation for asymptomatic screening for genital herpes.
False negative test results may occur if the test is performed late in the course of the illness or if the test sample is not appropriately acquired. Testing people for HSV when they are asymptomatic is not recommended due to the high false-positivity rate. A false positive test may cause relationship difficulties.
Genital herpes and pregnancy
Women who have genital herpes before pregnancy have a very low risk of transmitting herpes simplex virus to the baby during delivery. In the United States, 20-25% of pregnant women have genital herpes; however, fewer than 0.1% of babies born get neonatal herpes during delivery.
Per the U.S. Preventive Services Task Force, routine screening for pregnant women without a history of genital herpes is not recommended. Serologic (blood) antibody testing in asymptomatic patients without history has been shown to frequently have false positive and false negative test results which may lead to anxiety, labeling, or false reassurance with minimal improvements in health outcomes of reducing neonatal herpes transmission.
Pregnant women should notify their doctor if they show symptoms of genital herpes. At the time of delivery, women should be physically examined for signs of genital herpes. If a pregnant woman is symptomatic during delivery, a Cesarean section is the safest method of preventing contact and transmission of herpes simplex virus between the mother and the baby. Alternatively, some physicians use the drug acyclovir to treat pregnant women with genital herpes at 36 weeks until delivery to prevent the recurrence of symptoms and reduce the risk of transmission during delivery. Acyclovir is not approved for this purpose by the FDA; however, acyclovir's manufacturer has tracked pregnant women who have taken the drug during pregnancy, and there is no evidence that shows any risks for the infant.
Acyclovir may help reduce the frequency of symptomatic recurrence near term but may not definitively protect against transmission in all cases. This may be favorable particularly for women who prefer to have a vaginal delivery instead of cesarean section.
Treatment
There is no cure for the disease. Skin lesions disappear without treatment within a few weeks, but treatment accelerates the healing of lesions, reduces symptoms, and helps prevent or reduce recurrent outbreaks of the disease.
Antiviral medications provide clinical benefits to those who are symptomatic and is the primary means of management once infected. The main goal for the use of antiviral medications is to treat the first outbreak or to prevent genital herpes recurrences, improve quality of life, and help suppress the virus to sexual transmission to partners. Three FDA-approved antiviral medications have clinical benefits in controlling the signs and symptoms of genital herpes when used for first clinical symptoms and recurrent episodes or when used as daily suppressive therapy. These medications are acyclovir, valacyclovir, and famciclovir and have been shown to be safe with long-term use.
Acyclovir is an antiviral medication and reduces the pain and the number of lesions in the initial case of genital herpes. Furthermore, it decreases the frequency and severity of recurrent infections. It comes in capsules, tablets, and ointment. However, topical ointment with acyclovir is discouraged since it offers minimal clinical benefits.
Valacyclovir is a prodrug that is converted to acyclovir once in the body. It helps relieve the pain and discomfort and speeds healing of sores. It only comes in caplets and its advantage is that it has a longer duration of action than acyclovir.
Famciclovir is another antiviral drug that belongs to the same class. Famciclovir is a prodrug that is converted to penciclovir in the body. The latter is the one active against the viruses. It has a longer duration of action than acyclovir and it only comes in tablets.
First clinical episode of genital herpes
The first time an individual experiences genital herpes, they may have prolonged clinical illness with severe genital ulceration. Furthermore, for those who have mild clinical symptoms initially may experience severe recurrent infections later. Typical recommended regimens for first clinical episodes of genital herpes may be something like:
Acylovir 400 mg orally 3 times per day for 7–10 days or
Valacyclovir 1g orally 3 times per day for 7–10 days or
Famciclovir 1g orally 2 times per day for 7–10 days
A treatment longer than 10 days may be recommended if the genital ulcers have not fully healed.
Recurrent genital herpes
Most individuals who experience a symptomatic first episode of genital herpes will experience recurrence of genital lesions at some point in the future. Asymptomatic shedding can also occur where an individual may not have genital ulcerations present but still possibly transmit the virus to other partners. It is important for patients to have a discussion with their primary care doctor for options of receiving either episodic treatment or long-term suppressive therapies.
Suppressive therapy for recurrent genital herpes
Suppressive therapy has been shown effective in reducing recurrent genital herpes in as high as 80% which can tremendously help in improving quality of life since patients claim having minimal symptomatic episodes. Long-term use of anti-virals like acyclovir, valacyclovir, and famciclovir have been shown to be safe and effective. Furthermore, long-term treatment of genital herpes with valacyclovir daily has shown to decrease the rates of transmission. It is important for patients to continue suppressive therapy in conjunction to consistent condom use and sexual abstinence during recurrent episodes to decrease transmission as well.
Over-the-counter and non-drug treatments
To decrease symptoms during an outbreak of genital herpes, people can use an ice pack on the affected areas, take a warm bath, keep the genitals dry when not bathing, and take over-the-counter pain relief medication such as ibuprofen or acetaminophen.
Epidemiology
About 16 percent of Americans between the ages of 14 and 49 are infected with genital herpes, making it one of the most common sexually transmitted infections. More than 85% of those with HSV-2 are unaware of their infection. Approximately 776,000 people in the United States get new herpes infections every year.
Tests for herpes are not routinely included among STI screenings. Performers in the pornography industry are screened for HIV, chlamydia, and gonorrhea with an optional panel of tests for hepatitis B, hepatitis C and syphilis, but not herpes. Testing for herpes is controversial since the results are not always accurate or helpful. Most sex workers and performers will contract herpes at some point in their careers whether they use protection or not.
History
Early 20th century public health legislation in the United Kingdom required compulsory treatment for sexually transmitted infections but did not include herpes because it was not serious enough. As late as 1975, nursing textbooks did not include herpes as it was considered no worse than a common cold. After the development of acyclovir in the 1970s, the drug company Burroughs Wellcome launched an extensive marketing campaign that publicized the illness, including creating victim's support groups.
Research
There are efforts to develop a vaccine for active outbreaks of the virus—despite most cases being asymptomatic—but the results thus far have not been able to do so or eliminate transmission.
| Biology and health sciences | Viral diseases | Health |
5069516 | https://en.wikipedia.org/wiki/HIV/AIDS | HIV/AIDS | The human immunodeficiency virus (HIV) is a retrovirus that attacks the immune system. It is a preventable disease. It can be managed with treatment and become a manageable chronic health condition. While there is no cure or vaccine for HIV, antiretroviral treatment can slow the course of the disease and enable people living with HIV to lead long and healthy lives. An HIV-positive person on treatment can expect to live a normal life, and die with the virus, not of it. Effective treatment for HIV-positive people (people living with HIV) involves a life-long regimen of medicine to suppress the virus, making the viral load undetectable. Without treatment it can lead to a spectrum of conditions including acquired immunodeficiency syndrome (AIDS).
Treatment is recommended as soon as the diagnosis is made. An HIV-positive person who has an undetectable viral load as a result of long-term treatment has effectively no risk of transmitting HIV sexually. Campaigns by UNAIDS and organizations around the world have communicated this as Undetectable = Untransmittable. Without treatment the infection can interfere with the immune system, and eventually progress to AIDS, sometimes taking many years. Following initial infection an individual may not notice any symptoms, or may experience a brief period of influenza-like illness. During this period the person may not know that they are HIV-positive, yet they will be able to pass on the virus. Typically, this period is followed by a prolonged incubation period with no symptoms. Eventually the HIV infection increases the risk of developing other infections such as tuberculosis, as well as other opportunistic infections, and tumors which are rare in people who have normal immune function. The late stage is often also associated with unintended weight loss. Without treatment a person living with HIV can expect to live for 11 years. Early testing can show if treatment is needed to stop this progression and to prevent infecting others.
HIV is spread primarily by unprotected sex (including anal and vaginal sex), contaminated hypodermic needles or blood transfusions, and from mother to child during pregnancy, delivery, or breastfeeding. Some bodily fluids, such as saliva, sweat, and tears, do not transmit the virus. Oral sex has little risk of transmitting the virus. Ways to avoid catching HIV and preventing the spread include safe sex, treatment to prevent infection ("PrEP"), treatment to stop infection in someone who has been recently exposed ("PEP"), treating those who are infected, and needle exchange programs. Disease in a baby can often be prevented by giving both the mother and child antiretroviral medication.
Recognized worldwide in the early 1980s, HIV/AIDS has had a large impact on society, both as an illness and as a source of discrimination. The disease also has large economic impacts. There are many misconceptions about HIV/AIDS, such as the belief that it can be transmitted by casual non-sexual contact. The disease has become subject to many controversies involving religion, including the Catholic Church's position not to support condom use as prevention. It has attracted international medical and political attention as well as large-scale funding since it was identified in the 1980s.
HIV made the jump from other primates to humans in west-central Africa in the early-to-mid-20th century. AIDS was first recognized by the U.S. Centers for Disease Control and Prevention (CDC) in 1981 and its cause—HIV infection—was identified in the early part of the decade. Between the first time AIDS was readily identified through 2024, the disease is estimated to have caused at least 42.3 million deaths worldwide. In 2023, 630,000 people died from HIV-related causes, an estimated 1.3 million people acquired HIV and about 39.9 million people worldwide living with HIV, 65% of whom are in the World Health Organization (WHO) African Region. HIV/AIDS is considered a pandemic—a disease outbreak which is present over a large area and is actively spreading. The United States' National Institutes of Health (NIH) and the Gates Foundation have pledged $200 million focused on developing a global cure for AIDS.
Signs and symptoms
There are three main stages of HIV infection: acute infection, clinical latency, and AIDS.
First main stage: acute infection
The initial period following infection with HIV is called acute HIV, primary HIV or acute retroviral syndrome. Many individuals develop an illness like influenza, mononucleosis or glandular fever 2–4 weeks after exposure while others have no significant symptoms. Symptoms occur in 40–90% of cases and most commonly include fever, large tender lymph nodes, throat inflammation, a rash, headache, tiredness, and/or sores of the mouth and genitals. The rash, which occurs in 20–50% of cases, presents itself on the trunk and is maculopapular, classically. Some people also develop opportunistic infections at this stage. Gastrointestinal symptoms, such as vomiting or diarrhea may occur. Neurological symptoms of peripheral neuropathy or Guillain–Barré syndrome also occur. The duration of the symptoms varies, but is usually one or two weeks.
These symptoms are not often recognized as signs of HIV infection. Family doctors or hospitals can misdiagnose cases as one of the many common infectious diseases with similar symptoms. Someone with an unexplained fever who may have been recently exposed to HIV should consider testing to find out if they have been infected.
Second main stage: clinical latency
The initial symptoms are followed by a stage called clinical latency, asymptomatic HIV, or chronic HIV. Without treatment, this second stage of the natural history of HIV infection can last from about three years to over 20 years (on average, about eight years). While typically there are few or no symptoms at first, near the end of this stage many people experience fever, weight loss, gastrointestinal problems and muscle pains. Between 50% and 70% of people also develop persistent generalized lymphadenopathy, characterized by unexplained, non-painful enlargement of more than one group of lymph nodes (other than in the groin) for over three to six months.
Although most HIV-1 infected individuals have a detectable viral load and in the absence of treatment will eventually progress to AIDS, a small proportion (about 5%) retain high levels of CD4+ T cells (T helper cells) without antiretroviral therapy for more than five years. These individuals are classified as "HIV controllers" or long-term nonprogressors (LTNP). Another group consists of those who maintain a low or undetectable viral load without anti-retroviral treatment, known as "elite controllers" or "elite suppressors". They represent approximately 1 in 300 infected persons.
Third main stage: AIDS
Acquired immunodeficiency syndrome (AIDS) is defined as an HIV infection with either a CD4+ T cell count below 200 cells per μL or the occurrence of specific diseases associated with HIV infection. In the absence of specific treatment, around half of people infected with HIV develop AIDS within ten years. The most common initial conditions that alert to the presence of AIDS are pneumocystis pneumonia (40%), cachexia in the form of HIV wasting syndrome (20%), and esophageal candidiasis. Other common signs include recurrent respiratory tract infections.
Opportunistic infections may be caused by bacteria, viruses, fungi, and parasites that are normally controlled by the immune system. Which infections occur depends partly on what organisms are common in the person's environment. These infections may affect nearly every organ system.
People with AIDS have an increased risk of developing various viral-induced cancers, including Kaposi's sarcoma, Burkitt's lymphoma, primary central nervous system lymphoma, and cervical cancer. Kaposi's sarcoma is the most common cancer, occurring in 10% to 20% of people with HIV. The second-most common cancer is lymphoma, which is the cause of death of nearly 16% of people with AIDS and is the initial sign of AIDS in 3% to 4%. Both these cancers are associated with human herpesvirus 8 (HHV-8). Cervical cancer occurs more frequently in those with AIDS because of its association with human papillomavirus (HPV). Conjunctival cancer (of the layer that lines the inner part of eyelids and the white part of the eye) is also more common in those with HIV.
Additionally, people with AIDS frequently have systemic symptoms such as prolonged fevers, sweats (particularly at night), swollen lymph nodes, chills, weakness, and unintended weight loss. Diarrhea is another common symptom, present in about 90% of people with AIDS. They can also be affected by diverse psychiatric and neurological symptoms independent of opportunistic infections and cancers.
Transmission
HIV is spread by three main routes: sexual contact, significant exposure to infected body fluids or tissues, and from mother to child during pregnancy, delivery, or breastfeeding (known as vertical transmission). There is no risk of acquiring HIV if exposed to feces, nasal secretions, saliva, sputum, sweat, tears, urine, or vomit unless these are contaminated with blood. It is also possible to be co-infected by more than one strain of HIV—a condition known as HIV superinfection.
Sexual
The most frequent mode of transmission of HIV is through sexual contact with an infected person. However, an HIV-positive person who has an undetectable viral load as a result of long-term treatment has effectively no risk of transmitting HIV sexually, known as Undetectable = Untransmittable. The existence of functionally noncontagious HIV-positive people on antiretroviral therapy was controversially publicized in the 2008 Swiss Statement, and has since become accepted as medically sound.
Globally, the most common mode of HIV transmission is via sexual contacts between people of the opposite sex; however, the pattern of transmission varies among countries. , most HIV transmission in the United States occurred among men who had sex with men (82% of new HIV diagnoses among males aged 13 and older and 70% of total new diagnoses). In the US, gay and bisexual men aged 13 to 24 accounted for an estimated 92% of new HIV diagnoses among all men in their age group and 27% of new diagnoses among all gay and bisexual men.
With regard to unprotected heterosexual contacts, estimates of the risk of HIV transmission per sexual act appear to be four to ten times higher in low-income countries than in high-income countries. In low-income countries, the risk of female-to-male transmission is estimated as 0.38% per act, and of male-to-female transmission as 0.30% per act; the equivalent estimates for high-income countries are 0.04% per act for female-to-male transmission, and 0.08% per act for male-to-female transmission. The risk of transmission from anal intercourse is especially high, estimated as 1.4–1.7% per act in both heterosexual and homosexual contacts. While the risk of transmission from oral sex is relatively low, it is still present. The risk from receiving oral sex has been described as "nearly nil"; however, a few cases have been reported. The per-act risk is estimated at 0–0.04% for receptive oral intercourse. In settings involving prostitution in low-income countries, risk of female-to-male transmission has been estimated as 2.4% per act, and of male-to-female transmission as 0.05% per act.
Risk of transmission increases in the presence of many sexually transmitted infections and genital ulcers. Genital ulcers increase the risk approximately fivefold. Other sexually transmitted infections, such as gonorrhea, chlamydia, trichomoniasis, and bacterial vaginosis, are associated with somewhat smaller increases in risk of transmission.
The viral load of an infected person is an important risk factor in both sexual and mother-to-child transmission. During the first 2.5 months of an HIV infection, a person's infectiousness is twelve times higher due to the high viral load associated with acute HIV. If the person is in the late stages of infection, rates of transmission are approximately eightfold greater.
Commercial sex workers (including those in pornography) have an increased likelihood of contracting HIV. Rough sex can be a factor associated with an increased risk of transmission. Sexual assault is also believed to carry an increased risk of HIV transmission, as condoms are rarely worn, physical trauma to the vagina or rectum is likely, and there may be a greater risk of concurrent sexually transmitted infections.
Body fluids
The second-most frequent mode of HIV transmission is via blood and blood products. Blood-borne transmission can be through needle-sharing during intravenous drug use, needle-stick injury, transfusion of contaminated blood or blood product, or medical injections with unsterilized equipment. The risk from sharing a needle during drug injection is between 0.63% and 2.4% per act, with an average of 0.8%. The risk of acquiring HIV from a needle stick from an HIV-infected person is estimated as 0.3% (about 1 in 333) per act and the risk following mucous membrane exposure to infected blood as 0.09% (about 1 in 1000) per act. This risk may, however, be up to 5% if the introduced blood was from a person with a high viral load and the cut was deep. In the United States, intravenous drug users made up 12% of all new cases of HIV in 2009, and in some areas more than 80% of people who inject drugs are HIV-positive.
HIV is transmitted in about 90% of blood transfusions using infected blood. In developed countries the risk of acquiring HIV from a blood transfusion is extremely low (less than one in half a million) where improved donor selection and HIV screening is performed; for example, in the UK the risk was reported at one in five million in 2011 and in the United States it was one in 1.5 million in 2008. In low-income countries, only half of transfusions may be appropriately screened (as of 2008), and it is estimated that up to 15% of HIV infections in these areas come from transfusion of infected blood and blood products, representing between 5% and 10% of global infections. It is possible to acquire HIV from organ and tissue transplantation, although this is rare because of screening.
Unsafe medical injections play a role in HIV spread in sub-Saharan Africa. In 2007, between 12% and 17% of infections in this region were attributed to medical syringe use. The World Health Organization estimates the risk of transmission as a result of a medical injection in Africa at 1.2%. Risks are also associated with invasive procedures, assisted delivery, and dental care in this area of the world.
People giving or receiving tattoos, piercings, and scarification are theoretically at risk of infection but no confirmed cases have been documented. It is not possible for mosquitoes or other insects to transmit HIV.
Mother-to-child
HIV can be transmitted from mother to child during pregnancy, during delivery, or through breast milk, resulting in the baby also contracting HIV. As of 2008, vertical transmission accounted for about 90% of cases of HIV in children. In the absence of treatment, the risk of transmission before or during birth is around 20%, and in those who also breastfeed 35%. Treatment decreases this risk to less than 5%.
Antiretrovirals when taken by either the mother or the baby decrease the risk of transmission in those who do breastfeed. If blood contaminates food during pre-chewing it may pose a risk of transmission. If a woman is untreated, two years of breastfeeding results in an HIV/AIDS risk in her baby of about 17%. Due to the increased risk of death without breastfeeding in many areas in the developing world, the World Health Organization recommends either exclusive breastfeeding or the provision of safe formula. All women known to be HIV-positive should be taking lifelong antiretroviral therapy.
Virology
HIV is the cause of the spectrum of disease known as HIV/AIDS. HIV is a retrovirus that primarily infects components of the human immune system such as CD4+ T cells, macrophages and dendritic cells. It directly and indirectly destroys CD4+ T cells.
HIV is a member of the genus Lentivirus, part of the family Retroviridae. Lentiviruses share many morphological and biological characteristics. Many species of mammals are infected by lentiviruses, which are characteristically responsible for long-duration illnesses with a long incubation period. Lentiviruses are transmitted as single-stranded, positive-sense, enveloped RNA viruses. Upon entry into the target cell, the viral RNA genome is converted (reverse transcribed) into double-stranded DNA by a virally encoded reverse transcriptase that is transported along with the viral genome in the virus particle. The resulting viral DNA is then imported into the cell nucleus and integrated into the cellular DNA by a virally encoded integrase and host co-factors. Once integrated, the virus may become latent, allowing the virus and its host cell to avoid detection by the immune system. Alternatively, the virus may be transcribed, producing new RNA genomes and viral proteins that are packaged and released from the cell as new virus particles that begin the replication cycle anew.
HIV is now known to spread between CD4+ T cells by two parallel routes: cell-free spread and cell-to-cell spread, i.e. it employs hybrid spreading mechanisms. In the cell-free spread, virus particles bud from an infected T cell, enter the blood/extracellular fluid and then infect another T cell following a chance encounter. HIV can also disseminate by direct transmission from one cell to another by a process of cell-to-cell spread. The hybrid spreading mechanisms of HIV contribute to the virus' ongoing replication against antiretroviral therapies.
Two types of HIV have been characterized: HIV-1 and HIV-2. HIV-1 is the virus that was originally discovered (and initially referred to also as LAV or HTLV-III). It is more virulent, more infective, and is the cause of the majority of HIV infections globally. The lower infectivity of HIV-2 as compared with HIV-1 implies that fewer people exposed to HIV-2 will be infected per exposure. Because of its relatively poor capacity for transmission, HIV-2 is largely confined to West Africa.
Pathophysiology
After the virus enters the body, there is a period of rapid viral replication, leading to an abundance of virus in the peripheral blood. During primary infection, the level of HIV may reach several million virus particles per milliliter of blood. This response is accompanied by a marked drop in the number of circulating CD4+ T cells. The acute viremia is almost invariably associated with activation of CD8+ T cells, which kill HIV-infected cells, and subsequently with antibody production, or seroconversion. The CD8+ T cell response is thought to be important in controlling virus levels, which peak and then decline, as the CD4+ T cell counts recover. A good CD8+ T cell response has been linked to slower disease progression and a better prognosis, though it does not eliminate the virus.
Ultimately, HIV causes AIDS by depleting CD4+ T cells. This weakens the immune system and allows opportunistic infections. T cells are essential to the immune response and without them, the body cannot fight infections or kill cancerous cells. The mechanism of CD4+ T cell depletion differs in the acute and chronic phases. During the acute phase, HIV-induced cell lysis and killing of infected cells by CD8+ T cells accounts for CD4+ T cell depletion, although apoptosis may also be a factor. During the chronic phase, the consequences of generalized immune activation coupled with the gradual loss of the ability of the immune system to generate new T cells appear to account for the slow decline in CD4+ T cell numbers.
Although the symptoms of immune deficiency characteristic of AIDS do not appear for years after a person is infected, the bulk of CD4+ T cell loss occurs during the first weeks of infection, especially in the intestinal mucosa, which harbors the majority of the lymphocytes found in the body. The reason for the preferential loss of mucosal CD4+ T cells is that the majority of mucosal CD4+ T cells express the CCR5 protein which HIV uses as a co-receptor to gain access to the cells, whereas only a small fraction of CD4+ T cells in the bloodstream do so. A specific genetic change that alters the CCR5 protein when present in both chromosomes very effectively prevents HIV-1 infection.
HIV seeks out and destroys CCR5 expressing CD4+ T cells during acute infection. A vigorous immune response eventually controls the infection and initiates the clinically latent phase. CD4+ T cells in mucosal tissues remain particularly affected. Continuous HIV replication causes a state of generalized immune activation persisting throughout the chronic phase. Immune activation, which is reflected by the increased activation state of immune cells and release of pro-inflammatory cytokines, results from the activity of several HIV gene products and the immune response to ongoing HIV replication. It is also linked to the breakdown of the immune surveillance system of the gastrointestinal mucosal barrier caused by the depletion of mucosal CD4+ T cells during the acute phase of disease.
Diagnosis
HIV/AIDS is diagnosed via laboratory testing and then staged based on the presence of certain signs or symptoms. HIV screening is recommended by the United States Preventive Services Task Force for all people 15 years to 65 years of age, including all pregnant women. Additionally, testing is recommended for those at high risk, which includes anyone diagnosed with a sexually transmitted illness. In many areas of the world, a third of HIV carriers only discover they are infected at an advanced stage of the disease when AIDS or severe immunodeficiency has become apparent.
HIV testing
Most people infected with HIV develop seroconverted (antigen-specific) antibodies within three to twelve weeks after the initial infection. Diagnosis of primary HIV before seroconversion is done by measuring HIV-RNA or p24 antigen. Positive results obtained by antibody or PCR testing are confirmed either by a different antibody or by PCR.
Antibody tests in children younger than 18 months are typically inaccurate, due to the continued presence of maternal antibodies. Thus HIV infection can only be diagnosed by PCR testing for HIV RNA or DNA, or via testing for the p24 antigen. Much of the world lacks access to reliable PCR testing, and people in many places simply wait until either symptoms develop or the child is old enough for accurate antibody testing. In sub-Saharan Africa between 2007 and 2009, between 30% and 70% of the population were aware of their HIV status. In 2009, between 3.6% and 42% of men and women in sub-Saharan countries were tested; this represented a significant increase compared to previous years.
Classifications
Two main clinical staging systems are used to classify HIV and HIV-related disease for surveillance purposes: the WHO disease staging system for HIV infection and disease, and the CDC classification system for HIV infection. The CDC's classification system is more frequently adopted in developed countries. Since the WHO's staging system does not require laboratory tests, it is suited to the resource-restricted conditions encountered in developing countries, where it can also be used to help guide clinical management. Despite their differences, the two systems allow a comparison for statistical purposes.
The World Health Organization first proposed a definition for AIDS in 1986. Since then, the WHO classification has been updated and expanded several times, with the most recent version being published in 2007. The WHO system uses the following categories:
Primary HIV infection: May be either asymptomatic or associated with acute retroviral syndrome
Stage I: HIV infection is asymptomatic with a CD4+ T cell count (also known as CD4 count) greater than 500 per microlitre (μL or cubic mm) of blood. May include generalized lymph node enlargement.
Stage II: Mild symptoms, which may include minor mucocutaneous manifestations and recurrent upper respiratory tract infections. A CD4 count of less than 500/μL
Stage III: Advanced symptoms, which may include unexplained chronic diarrhea for longer than a month, severe bacterial infections including tuberculosis of the lung, and a CD4 count of less than 350/μL
Stage IV or AIDS: severe symptoms, which include toxoplasmosis of the brain, candidiasis of the esophagus, trachea, bronchi, or lungs, and Kaposi's sarcoma. A CD4 count of less than 200/μL
The U.S. Centers for Disease Control and Prevention also created a classification system for HIV, and updated it in 2008 and 2014. This system classifies HIV infections based on CD4 count and clinical symptoms, and describes the infection in five groups. In those greater than six years of age it is:
Stage 0: the time between a negative or indeterminate HIV test followed less than 180 days by a positive test
Stage 1: CD4 count ≥ 500 cells/μL and no AIDS-defining conditions
Stage 2: CD4 count 200 to 500 cells/μL and no AIDS-defining conditions
Stage 3: CD4 count ≤ 200 cells/μL or AIDS-defining conditions
Unknown: if insufficient information is available to make any of the above classifications.
For surveillance purposes, the AIDS diagnosis still stands even if, after treatment, the CD4+ T cell count rises to above 200 per μL of blood or other AIDS-defining illnesses are cured.
Prevention
Sexual contact
Consistent condom use reduces the risk of HIV transmission by approximately 80% over the long term. When condoms are used consistently by a couple in which one person is infected, the rate of HIV infection is less than 1% per year. There is some evidence to suggest that female condoms may provide an equivalent level of protection. Application of a vaginal gel containing tenofovir (a reverse transcriptase inhibitor) immediately before sex seems to reduce infection rates by approximately 40% among African women. By contrast, use of the spermicide nonoxynol-9 may increase the risk of transmission due to its tendency to cause vaginal and rectal irritation.
Circumcision in sub-Saharan Africa "reduces the acquisition of HIV by heterosexual men by between 38% and 66% over 24 months". Owing to these studies, both the World Health Organization and UNAIDS recommended male circumcision in 2007 as a method of preventing female-to-male HIV transmission in areas with high rates of HIV. However, whether it protects against male-to-female transmission is disputed, and whether it is of benefit in developed countries and among men who have sex with men is undetermined.
Programs encouraging sexual abstinence do not appear to affect subsequent HIV risk. Evidence of any benefit from peer education is equally poor. Comprehensive sexual education provided at school may decrease high-risk behavior. A substantial minority of young people continues to engage in high-risk practices despite knowing about HIV/AIDS, underestimating their own risk of becoming infected with HIV. Voluntary counseling and testing people for HIV does not affect risky behavior in those who test negative but does increase condom use in those who test positive. Enhanced family planning services appear to increase the likelihood of women with HIV using contraception, compared to basic services. It is not known whether treating other sexually transmitted infections is effective in preventing HIV.
Pre-exposure
Antiretroviral treatment among people with HIV whose CD4 count ≤ 550 cells/μL is a very effective way to prevent HIV infection of their partner (a strategy known as treatment as prevention, or TASP). TASP is associated with a 10- to 20-fold reduction in transmission risk. Pre-exposure prophylaxis for HIV ("PrEP") with a daily dose of the medications tenofovir, with or without emtricitabine, is effective in people at high risk including men who have sex with men, couples where one is HIV-positive, and young heterosexuals in Africa. It may also be effective in intravenous drug users, with a study finding a decrease in risk of 0.7 to 0.4 per 100 person years. The USPSTF, in 2019, recommended PrEP in those who are at high risk.
Universal precautions within the health care environment are believed to be effective in decreasing the risk of HIV. Intravenous drug use is an important risk factor, and harm reduction strategies such as needle-exchange programs and opioid substitution therapy appear effective in decreasing this risk.
Post-exposure
A course of antiretrovirals administered within 48 to 72 hours after exposure to HIV-positive blood or genital secretions is referred to as post-exposure prophylaxis (PEP). The use of the single agent zidovudine reduces the risk of an HIV infection five-fold following a needle-stick injury. , the prevention regimen recommended in the United States consists of three medications—tenofovir, emtricitabine and raltegravir—as this may reduce the risk further.
PEP treatment is recommended after a sexual assault when the perpetrator is known to be HIV-positive, but is controversial when their HIV status is unknown. The duration of treatment is usually four weeks and is frequently associated with adverse effects—where zidovudine is used, about 70% of cases result in adverse effects such as nausea (24%), fatigue (22%), emotional distress (13%) and headaches (9%).
Mother-to-child
Programs to prevent the vertical transmission of HIV (from mothers to children) can reduce rates of transmission by 92–99%. This primarily involves the use of a combination of antiviral medications during pregnancy and after birth in the infant, and potentially includes bottle feeding rather than breastfeeding. If replacement feeding is acceptable, feasible, affordable, sustainable and safe, mothers should avoid breastfeeding their infants; however, exclusive breastfeeding is recommended during the first months of life if this is not the case. If exclusive breastfeeding is carried out, the provision of extended antiretroviral prophylaxis to the infant decreases the risk of transmission. In 2015, Cuba became the first country in the world to eradicate mother-to-child transmission of HIV.
Vaccination
Currently there is no licensed vaccine for HIV or AIDS. The most effective vaccine trial to date, RV 144, was published in 2009; it found a partial reduction in the risk of transmission of roughly 30%, stimulating some hope in the research community of developing a truly effective vaccine.
Treatment
HIV/AIDS is a terminal illness, as there is currently no cure, nor an effective HIV vaccine. Treatment consists of highly active antiretroviral therapy (ART), which slows progression of the disease. As of 2022, 39 million people globally were living with HIV, and 29.8 million people were accessing ART. Treatment also includes preventive and active treatment of opportunistic infections. , four people have been successfully cleared of HIV. Rapid initiation of antiretroviral therapy within one week of diagnosis appear to improve treatment outcomes in low and medium-income settings and is recommend for newly diagnosed HIV patients.
Antiviral therapy
Current ART options are combinations (or "cocktails") consisting of at least three medications belonging to at least two types, or "classes", of antiretroviral agents. There are eight classes of antiretroviral agents (ARVs), and over 30 individual drugs: nucleoside/nucleotide reverse transcriptase inhibitors (NRTIs), non-nucleoside reverse transcriptase, inhibitors (NNRTIs), protease inhibitors (PIs), integrase strand transfer inhibitors (INSTIs), a fusion inhibitor, a CCR5 antagonist, a CD4 T lymphocyte (CD4) post-attachment inhibitor, and a gp120 attachment inhibitor. There are also two drugs, ritonavir (RTV) and cobicistat (COBI) which can be used as pharmacokinetic (PK) enhancers (or boosters) to improve the PK profiles of PIs and the INSTI elvitegravir (EVG). Depending on the guidelines being followed, initial treatment generally consists of two nucleoside reverse transcriptase inhibitors along with a third ARV, either an integrase strand transfer inhibitor (INSTI), a non-nucleoside reverse transcriptase inhibitor (NNRTI), or a protease inhibitor with a pharmacokinetic enhancer (also known as a booster).
The World Health Organization and the United States recommend antiretrovirals in people of all ages (including pregnant women) as soon as the diagnosis is made, regardless of CD4 count. Once treatment is begun, it is recommended that it is continued without breaks or "holidays". Many people are diagnosed only after treatment ideally should have begun. The desired outcome of treatment is a long-term plasma HIV-RNA count below 50 copies/mL. Levels to determine if treatment is effective are initially recommended after four weeks and once levels fall below 50 copies/mL checks every three to six months are typically adequate. Inadequate control is deemed to be greater than 400 copies/mL. Based on these criteria treatment is effective in more than 95% of people during the first year.
Benefits of treatment include a decreased risk of progression to AIDS and a decreased risk of death. In the developing world, treatment also improves physical and mental health. With treatment, there is a 70% reduced risk of acquiring tuberculosis. Additional benefits include a decreased risk of transmission of the disease to sexual partners and a decrease in mother-to-child transmission. The effectiveness of treatment depends to a large part on compliance. Reasons for non-adherence to treatment include poor access to medical care, inadequate social supports, mental illness and drug abuse. The complexity of treatment regimens (due to pill numbers and dosing frequency) and adverse effects may reduce adherence. Even though cost is an important issue with some medications, 47% of those who needed them were taking them in low- and middle-income countries , and the rate of adherence is similar in low-income and high-income countries.
Specific adverse events are related to the antiretroviral agent taken. Some relatively common adverse events include: lipodystrophy syndrome, dyslipidemia, and diabetes mellitus, especially with protease inhibitors. Other common symptoms include diarrhea, and an increased risk of cardiovascular disease. Newer recommended treatments are associated with fewer adverse effects. Certain medications may be associated with birth defects and therefore may be unsuitable for women hoping to have children.
Treatment recommendations for children are somewhat different from those for adults. The World Health Organization recommends treating all children less than five years of age; children above five are treated like adults. The United States guidelines recommend treating all children less than 12 months of age and all those with HIV RNA counts greater than 100,000 copies/mL between one year and five years of age.
The European Medicines Agency (EMA) has recommended the granting of marketing authorizations for two new antiretroviral (ARV) medicines, rilpivirine (Rekambys) and cabotegravir (Vocabria), to be used together for the treatment of people with human immunodeficiency virus type 1 (HIV-1) infection. The two medicines are the first ARVs that come in a long-acting injectable formulation. This means that instead of daily pills, people receive intramuscular injections monthly or every two months.
The combination of Rekambys and Vocabria injection is intended for maintenance treatment of adults who have undetectable HIV levels in the blood (viral load less than 50 copies/mL) with their current ARV treatment, and when the virus has not developed resistance to a certain class of anti-HIV medicines called non-nucleoside reverse transcriptase inhibitors (NNRTIs) and integrase strand transfer inhibitors (INIs).
Cabotegravir combined with rilpivirine (Cabenuva) is a complete regimen for the treatment of human immunodeficiency virus type 1 (HIV-1) infection in adults to replace a current antiretroviral regimen in those who are virologically suppressed on a stable antiretroviral regimen with no history of treatment failure and with no known or suspected resistance to either cabotegravir or rilpivirine.
Opportunistic infections
Measures to prevent opportunistic infections are effective in many people with HIV/AIDS. In addition to improving current disease, treatment with antiretrovirals reduces the risk of developing additional opportunistic infections.
Adults and adolescents who are living with HIV (even on anti-retroviral therapy) with no evidence of active tuberculosis in settings with high tuberculosis burden should receive isoniazid preventive therapy (IPT); the tuberculin skin test can be used to help decide if IPT is needed. Children with HIV may benefit from screening for tuberculosis. Vaccination against hepatitis A and B is advised for all people at risk of HIV before they become infected; however, it may also be given after infection.
Trimethoprim/sulfamethoxazole prophylaxis between four and six weeks of age, and ceasing breastfeeding of infants born to HIV-positive mothers, is recommended in resource-limited settings. It is also recommended to prevent PCP when a person's CD4 count is below 200 cells/uL and in those who have or have previously had PCP. People with substantial immunosuppression are also advised to receive prophylactic therapy for toxoplasmosis and MAC. Appropriate preventive measures reduced the rate of these infections by 50% between 1992 and 1997. Influenza vaccination and pneumococcal polysaccharide vaccine are often recommended in people with HIV/AIDS with some evidence of benefit.
Diet
The World Health Organization (WHO) has issued recommendations regarding nutrient requirements in HIV/AIDS. A generally healthy diet is promoted. Dietary intake of micronutrients at RDA levels by HIV-infected adults is recommended by the WHO; higher intake of vitamin A, zinc, and iron can produce adverse effects in HIV-positive adults, and is not recommended unless there is documented deficiency. Dietary supplementation for people who are infected with HIV and who have inadequate nutrition or dietary deficiencies may strengthen their immune systems or help them recover from infections; however, evidence indicating an overall benefit in morbidity or reduction in mortality is not consistent.
People with HIV/AIDS are up to four times more likely to develop type 2 diabetes than those who are not tested positive with the virus.
Evidence for supplementation with selenium is mixed with some tentative evidence of benefit. For pregnant and lactating women with HIV, multivitamin supplement improves outcomes for both mothers and children. If the pregnant or lactating mother has been advised to take anti-retroviral medication to prevent mother-to-child HIV transmission, multivitamin supplements should not replace these treatments. There is some evidence that vitamin A supplementation in children with an HIV infection reduces mortality and improves growth.
Alternative medicine
In the US, approximately 60% of people with HIV use various forms of complementary or alternative medicine, whose effectiveness has not been established. There is not enough evidence to support the use of herbal medicines. There is insufficient evidence to recommend or support the use of medical cannabis to try to increase appetite or weight gain.
Prognosis
HIV/AIDS has become a chronic rather than an acutely fatal disease in many areas of the world. Prognosis varies between people, and both the CD4 count and viral load are useful for predicted outcomes. Without treatment, average survival time after infection with HIV is estimated to be 9 to 11 years, depending on the HIV subtype. After the diagnosis of AIDS, if treatment is not available, survival ranges between 6 and 19 months. ART and appropriate prevention of opportunistic infections reduces the death rate by 80%, and raises the life expectancy for a newly diagnosed young adult to 20–50 years. This is between two thirds and nearly that of the general population. If treatment is started late in the infection, prognosis is not as good: for example, if treatment is begun following the diagnosis of AIDS, life expectancy is ~10–40 years. Half of infants born with HIV die before two years of age without treatment.
The primary causes of death from HIV/AIDS are opportunistic infections and cancer, both of which are frequently the result of the progressive failure of the immune system. Risk of cancer appears to increase once the CD4 count is below 500/μL. The rate of clinical disease progression varies widely between individuals and has been shown to be affected by a number of factors such as a person's susceptibility and immune function; their access to health care, the presence of co-infections; and the particular strain (or strains) of the virus involved.
Tuberculosis co-infection is one of the leading causes of sickness and death in those with HIV/AIDS being present in a third of all HIV-infected people and causing 25% of HIV-related deaths. HIV is also one of the most important risk factors for tuberculosis. Hepatitis C is another very common co-infection where each disease increases the progression of the other. The two most common cancers associated with HIV/AIDS are Kaposi's sarcoma and AIDS-related non-Hodgkin's lymphoma. Other cancers that are more frequent include anal cancer, Burkitt's lymphoma, primary central nervous system lymphoma, and cervical cancer.
Even with anti-retroviral treatment, over the long term HIV-infected people may experience neurocognitive disorders, osteoporosis, neuropathy, cancers, nephropathy, and cardiovascular disease. Some conditions, such as lipodystrophy, may be caused both by HIV and its treatment.
Epidemiology
HIV/AIDS is considered a global pandemic. , approximately 39.0 million people worldwide are living with HIV, the number of new infections that year being about 1.3 million. This is down from 2.1 million new infections in 2010. Among new infections, 46% are in women and are children globally. There were 630,000 AIDS related deaths in 2022, down from a peak of 2 million in 2005. The World Health Organization has reported that deaths from HIV and AIDS have "fallen by 61%, moving from the world’s seventh leading cause of death in 2000 to the twenty-first in 2021."
Among persons living with HIV (PLWH), the largest proportion reside in eastern and southern Africa (20.6 million, 54.6%). This region also had the highest rate of adult and child deaths due to AIDS in 2020 (310,000, 46.6%). Sub-Saharan African adolescent girls and young women (aged 15–24 years) account for 77% of new infections among this age-range globally. Here, in contrast to other regions, adolescent girls and young women are three times more likely to acquire HIV than age-matched males. Despite these statistics, overall, new HIV infections and AIDS-related deaths have substantially decreased in this region since 2010.
Eastern Europe and central Asia has observed a 43% increase in new HIV infections and 32% increase in AIDS-related deaths since 2010, the highest of all global regions. These infections are predominantly distributed in persons who inject drugs, with gay men and other men who have sex with men or persons who engage in transaction sex the second and third populations most impacted in this region.
At the end of 2019, United States indicated that approximately 1.2 million people aged ≥13 years were living with HIV, resulting in about 18,500 deaths in 2020. There were 34,800 estimated new infections in the US in 2019, 53% of which were in the southern region of the country. In addition to geographic location, significant disparities in HIV incidence exist among men, Black or Hispanic populations, and men who reported male-to-male sexual contact. The US Centers for Disease Control and Prevention estimated that in that year, 158,500 people or 13% of infected Americans were unaware of their infection.
In the United Kingdom , there were approximately 101,200 cases which resulted in 594 deaths. In Canada as of 2008, there were about 65,000 cases causing 53 deaths. Between the first recognition of AIDS (in 1981) and 2009, it has led to nearly 30 million deaths. Rates of HIV are lowest in North Africa and the Middle East (0.1% or less), East Asia (0.1%), and Western and Central Europe (0.2%). The worst-affected European countries, in 2009 and 2012 estimates, are Russia, Ukraine, Latvia, Moldova, Portugal and Belarus, in decreasing order of prevalence.
Groups at higher risk of acquiring HIV include persons who engage in transactional sex, gay men and other men who have sex with men, persons who inject drugs, transgender persons, and those who are incarcerated or detained.
History
Discovery
The first news story on the disease appeared on May 18, 1981, in the gay newspaper New York Native. AIDS was first clinically reported on June 5, 1981, with five cases in the United States. The initial cases were a cluster of injecting drug users and gay men with no known cause of impaired immunity who showed symptoms of Pneumocystis carinii pneumonia (PCP), a rare opportunistic infection that was known to occur in people with very compromised immune systems. Soon thereafter, a large number of homosexual men developed a generally rare skin cancer called Kaposi's sarcoma (KS). Many more cases of PCP and KS emerged, alerting U.S. Centers for Disease Control and Prevention (CDC) and a CDC task force was formed to monitor the outbreak.
In the early days, the CDC did not have an official name for the disease, often referring to it by way of diseases associated with it, such as lymphadenopathy, the disease after which the discoverers of HIV originally named the virus. They also used Kaposi's sarcoma and opportunistic infections, the name by which a task force had been set up in 1981. At one point the CDC referred to it as the "4H disease", as the syndrome seemed to affect heroin users, homosexuals, hemophiliacs, and Haitians. The term GRID, which stood for gay-related immune deficiency, had also been coined. However, after determining that AIDS was not isolated to the gay community, it was realized that the term GRID was misleading, and the term AIDS was introduced at a meeting in July 1982. By September 1982 the CDC started referring to the disease as AIDS.
In 1983, two separate research groups led by Robert Gallo and Luc Montagnier declared that a novel retrovirus may have been infecting people with AIDS, and published their findings in the same issue of the journal Science. Gallo claimed a virus which his group had isolated from a person with AIDS was strikingly similar in shape to other human T-lymphotropic viruses (HTLVs) that his group had been the first to isolate. Gallo's group called their newly isolated virus HTLV-III. At the same time, Montagnier's group isolated a virus from a person presenting with swelling of the lymph nodes of the neck and physical weakness, two characteristic symptoms of AIDS. Contradicting the report from Gallo's group, Montagnier and his colleagues showed that core proteins of this virus were immunologically different from those of HTLV-I. Montagnier's group named their isolated virus lymphadenopathy-associated virus (LAV). As these two viruses turned out to be the same, in 1986, LAV and HTLV-III were renamed HIV.
Origins
The origin of HIV / AIDS and the circumstances that led to its emergence remain unsolved.
Both HIV-1 and HIV-2 are believed to have originated in non-human primates in West-central Africa and were transferred to humans in the early 20th century. HIV-1 appears to have originated in southern Cameroon through the evolution of SIV(cpz), a simian immunodeficiency virus (SIV) that infects wild chimpanzees (HIV-1 descends from the SIVcpz endemic in the chimpanzee subspecies Pan troglodytes troglodytes). The closest relative of HIV-2 is SIV (smm), a virus of the sooty mangabey (Cercocebus atys atys), an Old World monkey living in coastal West Africa (from southern Senegal to western Ivory Coast). New World monkeys such as the owl monkey are resistant to HIV-1 infection, possibly because of a genomic fusion of two viral resistance genes. HIV-1 is thought to have jumped the species barrier on at least three separate occasions, giving rise to the three groups of the virus, M, N, and O.
There is evidence that humans who participate in bushmeat activities, either as hunters or as bushmeat vendors, commonly acquire SIV. However, SIV is a weak virus which is typically suppressed by the human immune system within weeks of infection. It is thought that several transmissions of the virus from individual to individual in quick succession are necessary to allow it enough time to mutate into HIV. Furthermore, due to its relatively low person-to-person transmission rate, SIV can only spread throughout the population in the presence of one or more high-risk transmission channels, which are thought to have been absent in Africa before the 20th century.
Specific proposed high-risk transmission channels, allowing the virus to adapt to humans and spread throughout society, depend on the proposed timing of the animal-to-human crossing. Genetic studies of the virus suggest that the most recent common ancestor of the HIV-1 M group dates back to 1910. Proponents of this dating link the HIV epidemic with the emergence of colonialism and growth of large colonial African cities, leading to social changes, including a higher degree of sexual promiscuity, the spread of prostitution, and the accompanying high frequency of genital ulcer diseases (such as syphilis) in nascent colonial cities. While transmission rates of HIV during vaginal intercourse are low under regular circumstances, they are increased manyfold if one of the partners has a sexually transmitted infection causing genital ulcers. Early 1900s colonial cities were notable for their high prevalence of prostitution and genital ulcers, to the degree that, as of 1928, as many as 45% of female residents of eastern Kinshasa were thought to have been prostitutes, and, as of 1933, around 15% of all residents of the same city had syphilis.
An alternative view holds that unsafe medical practices in Africa after World War II, such as unsterile reuse of single-use syringes during mass vaccination, antibiotic and anti-malaria treatment campaigns, were the initial vector that allowed the virus to adapt to humans and spread.
The earliest well-documented case of HIV in a human dates back to 1959 in the Congo. The virus may have been present in the U.S. as early as the mid-to-late 1950s. A 16-year-old male named Robert Rayford presented with symptoms in 1966 and died in 1969. In the 1970s, there were cases of people getting parasites and becoming sick with what was then called "gay bowel disease" but is now suspected to have been AIDS.
The earliest retrospectively described case of AIDS is believed to have been in Norway beginning in 1966, that of Arvid Noe. In July 1960, in the wake of Congo's independence, the United Nations recruited Francophone experts and technicians from all over the world to assist in filling administrative gaps left by Belgium, who did not leave behind an African elite to run the country. By 1962, Haitians made up the second-largest group of well-educated experts (out of the 48 national groups recruited), that totaled around 4,500 in the country. Dr. Jacques Pépin, a Canadian author of The Origins of AIDS, stipulates that Haiti was one of HIV's entry points to the U.S. and that a Haitian may have carried HIV back across the Atlantic in the 1960s. Although there was known to have been at least one case of AIDS in the U.S. from 1966, the vast majority of infections occurring outside sub-Saharan Africa (including the U.S.) can be traced back to a single unknown individual who became infected with HIV in Haiti and brought the infection to the U.S. at some time around 1969. The epidemic rapidly spread among high-risk groups (initially, sexually promiscuous men who have sex with men). By 1978, the prevalence of HIV-1 among gay male residents of New York City and San Francisco was estimated at 5%, suggesting that several thousand individuals in the country had been infected.
Society and culture
Stigma
AIDS stigma exists around the world in a variety of ways, including ostracism, rejection, discrimination and avoidance of HIV-infected people; compulsory HIV testing without prior consent or protection of confidentiality; violence against HIV-infected individuals or people who are perceived to be infected with HIV; and the quarantine of HIV-infected individuals. Stigma-related violence or the fear of violence prevents many people from seeking HIV testing, returning for their results, or securing treatment, possibly turning what could be a manageable chronic illness into a death sentence and perpetuating the spread of HIV.
AIDS stigma has been further divided into the following three categories:
Instrumental AIDS stigma—a reflection of the fear and apprehension that are likely to be associated with any deadly and transmissible illness.
Symbolic AIDS stigma—the use of HIV/AIDS to express attitudes toward the social groups or lifestyles perceived to be associated with the disease.
Courtesy AIDS stigma—stigmatization of people connected to the issue of HIV/AIDS or HIV-positive people.
Often, AIDS stigma is expressed in conjunction with one or more other stigmas, particularly those associated with homosexuality, bisexuality, promiscuity, prostitution, and intravenous drug use.
In many developed countries, there is an association between AIDS and homosexuality or bisexuality, and this association is correlated with higher levels of sexual prejudice, such as anti-homosexual or anti-bisexual attitudes. There is also a perceived association between AIDS and all male-male sexual behavior, including sex between uninfected men. However, the dominant mode of spread worldwide for HIV remains heterosexual transmission.
The NAMES Project AIDS Memorial Quilt was conceived in 1985 to celebrate the lives of those who had died of AIDS when stigma prevented many from receiving funerals. It is now cared for by the National AIDS Memorial in San Francisco.
In 2003, as part of an overall reform of marriage and population legislation, it became legal for those diagnosed with AIDS to marry in China.
Between 2004 and 2020, Somen Debnath has travelled the world by bicycle promoting HIV / AIDS awareness.
In 2013, the U.S. National Library of Medicine developed a traveling exhibition titled Surviving and Thriving: AIDS, Politics, and Culture; this covered medical research, the U.S. government's response, and personal stories from people with AIDS, caregivers, and activists.
Stigma has proved an obstacle to the update of PrEP. Within the MSM community, the greatest barrier to PrEP use has been the stigma surrounding HIV and gay men. Gay men on PrEP have experienced "slut-shaming". Numerous other barriers have been identified, including lack of quality LGBTQ care, cost, and adherence to medication use.
Economic impact
HIV/AIDS affects the economics of both individuals and countries. The gross domestic product of the most affected countries has decreased due to the lack of human capital. Without proper nutrition, health care and medicine, large numbers of people die from AIDS-related complications. Before death they will not only be unable to work, but will also require significant medical care. It is estimated that as of 2007 there were 12 million AIDS orphans. Many are cared for by elderly grandparents.
Returning to work after beginning treatment for HIV/AIDS is difficult, and affected people often work less than the average worker. Unemployment in people with HIV/AIDS also is associated with suicidal ideation, memory problems, and social isolation. Employment increases self-esteem, sense of dignity, confidence, and quality of life for people with HIV/AIDS. Anti-retroviral treatment may help people with HIV/AIDS work more, and may increase the chance that a person with HIV/AIDS will be employed (low-quality evidence).
By affecting mainly young adults, AIDS reduces the taxable population, in turn reducing the resources available for public expenditures such as education and health services not related to AIDS, resulting in increasing pressure on the state's finances and slower growth of the economy. This causes a slower growth of the tax base, an effect that is reinforced if there are growing expenditures on treating the sick, training (to replace sick workers), sick pay, and caring for AIDS orphans. This is especially true if the sharp increase in adult mortality shifts the responsibility from the family to the government in caring for these orphans.
At the household level, AIDS causes both loss of income and increased spending on healthcare. A study in Côte d'Ivoire showed that households having a person with HIV/AIDS spent twice as much on medical expenses as other households. This additional expenditure also leaves less income to spend on education and other personal or family investment.
Religion and AIDS
The topic of religion and AIDS has become highly controversial, primarily because some religious authorities have publicly declared their opposition to the use of condoms. The religious approach to prevent the spread of AIDS, according to a report by American health expert Matthew Hanley titled The Catholic Church and the Global AIDS Crisis, argues that cultural changes are needed, including a re-emphasis on fidelity within marriage and sexual abstinence outside of it.
Some religious organizations have claimed that prayer can cure HIV/AIDS. In 2011, the BBC reported that some churches in London were claiming that prayer would cure AIDS, and the Hackney-based Centre for the Study of Sexual Health and HIV reported that several people stopped taking their medication, sometimes on the direct advice of their pastor, leading to many deaths. The Synagogue Church Of All Nations advertised an "anointing water" to promote God's healing, although the group denies advising people to stop taking medication.
Media portrayal
One of the first high-profile cases of AIDS was the American gay actor Rock Hudson. He had been diagnosed during 1984, announced that he had had the virus on July 25, 1985, and died a few months later on October 2, 1985. Another notable British casualty of AIDS that year was Nicholas Eden, a gay politician and son of former prime minister Anthony Eden. On November 24, 1991, British rock star Freddie Mercury died from an AIDS-related illness, having revealed the diagnosis only on the previous day.
One of the first high-profile heterosexual cases of the virus was American tennis player Arthur Ashe. He was diagnosed as HIV-positive on August 31, 1988, having contracted the virus from blood transfusions during heart surgery earlier in the 1980s. Further tests within 24 hours of the initial diagnosis revealed that Ashe had AIDS, but he did not tell the public about his diagnosis until April 1992. He died as a result on February 6, 1993, aged 49.
Therese Frare's photograph of gay activist David Kirby, as he lay dying from AIDS while surrounded by family, was taken in April 1990. Life magazine said the photo became the one image "most powerfully identified with the HIV/AIDS epidemic." The photo was displayed in Life, was the winner of the World Press Photo, and acquired worldwide notoriety after being used in a United Colors of Benetton advertising campaign in 1992.
Many famous artists and AIDS activists such as Larry Kramer, Diamanda Galás and Rosa von Praunheim campaign for AIDS education and the rights of those affected. These artists worked with various media formats.
Criminal transmission
Criminal transmission of HIV is the intentional or reckless infection of a person with the human immunodeficiency virus (HIV). Some countries or jurisdictions, including some areas of the United States, have laws that criminalize HIV transmission or exposure. Others may charge the accused under laws enacted before the HIV pandemic.
In 1996, Ugandan-born Canadian Johnson Aziga was diagnosed with HIV; he subsequently had unprotected sex with eleven women without disclosing his diagnosis. By 2003, seven had contracted HIV; two died from complications related to AIDS. Aziga was convicted of first-degree murder and sentenced to life imprisonment.
Misconceptions
There are many misconceptions about HIV and AIDS. Three misconceptions are that AIDS can spread through casual contact, that sexual intercourse with a virgin will cure AIDS, and that HIV can infect only gay men and drug users. In 2014, some among the British public wrongly thought one could get HIV from kissing (16%), sharing a glass (5%), spitting (16%), a public toilet seat (4%), and coughing or sneezing (5%). Other misconceptions are that any act of anal intercourse between two uninfected gay men can lead to HIV infection, and that open discussion of HIV and homosexuality in schools will lead to increased rates of AIDS.
A small group of individuals continue to dispute the connection between HIV and AIDS, the existence of HIV itself, or the validity of HIV testing and treatment methods. These claims, known as AIDS denialism, have been examined and rejected by the scientific community. However, they have had a significant political impact, particularly in South Africa, where the government's official embrace of AIDS denialism (1999–2005) was responsible for its ineffective response to that country's AIDS epidemic, and has been blamed for hundreds of thousands of avoidable deaths and HIV infections.
Several discredited conspiracy theories have held that HIV was created by scientists, either inadvertently or deliberately. Operation INFEKTION was a worldwide Soviet active measures operation to spread the claim that the United States had created HIV/AIDS. Surveys show that a significant number of people believed—and continue to believe—in such claims.
Research
HIV/AIDS research includes all medical research which attempts to prevent, treat, or cure HIV/AIDS, along with fundamental research about the nature of HIV as an infectious agent, and about AIDS as the disease caused by HIV.
Many governments and research institutions participate in HIV/AIDS research. This research includes behavioral health interventions such as sex education, and drug development, such as research into microbicides for sexually transmitted diseases, HIV vaccines, and antiretroviral drugs. Other medical research areas include the topics of pre-exposure prophylaxis, post-exposure prophylaxis, and circumcision and HIV. Public health officials, researchers, and programs can gain a more comprehensive picture of the barriers they face, and the efficacy of current approaches to HIV treatment and prevention, by tracking standard HIV indicators. Use of common indicators is an increasing focus of development organizations and researchers.
| Biology and health sciences | Illness and injury | null |
5069898 | https://en.wikipedia.org/wiki/Dynamic%20demand%20%28electric%20power%29 | Dynamic demand (electric power) | Dynamic Demand is the name of a semi-passive technology to support demand response by adjusting the load demand on an electrical power grid. It is also the name of an independent not-for-profit organization in the UK supported by a charitable grant from the Esmée Fairbairn Foundation, dedicated to promoting this technology. The concept is that by monitoring the frequency of the power grid, as well as their own controls, intermittent domestic and industrial loads switch themselves on/off at optimal moments to balance the overall grid load with generation, reducing critical power mismatches. As this switching would only advance or delay the appliance operating cycle by a few seconds, it would be unnoticeable to the end user. This is the foundation of dynamic demand control. In the United States, in 1982, a (now-lapsed) patent for this idea was issued to power systems engineer Fred Schweppe. Other patents have been issued based on this idea.
Dynamic demand is similar to demand response mechanisms to manage domestic and industrial consumption of electricity in response to supply conditions, for example, having electricity customers reduce their consumption at critical times or in response to prices. The difference is that dynamic demand devices passively shut off when stress in the grid is sensed, whereas demand response mechanisms respond to transmitted requests to shut off,
The need for spinning reserve
The power utilities are able to predict to a reasonable accuracy (generally to within one or two percent) the demand pattern throughout any particular day. This means that the free market in electricity is able to schedule just enough base load in advance. Any remaining imbalance would then be due either to inaccuracies in the prediction, or unscheduled changes in supply (such as a power station fault) and/or demand. Such imbalances are removed by requesting generators to operate in so called frequency response mode (also called frequency control mode), altering their output continuously to keep the frequency near the required value.
The grid frequency is a system-wide indicator of overall power imbalance. For example, it will drop if there is too much demand because generators will start to slow down slightly. A generator in frequency-response mode will, under nominal conditions, run at reduced output in order to maintain a buffer of spare capacity. It will then continually alter its output on a second-to-second basis to the needs of the grid with droop speed control.
This spinning reserve is a significant expense to the power utilities as often fuel must be burned or potential power sales lost to maintain it. The kind of generation used for fast response is usually fossil fuel powered which produces emissions of between 0.48 and 1.3 tonnes of CO2 equivalent for every megawatt hour (MWh) generated. Thus, a significant environmental burden, in the form of increased greenhouse gas emissions, is associated with this imbalance.
Local load control
In principle, any appliance that operates to a duty cycle (such as industrial or domestic air conditioners, water heaters, heat pumps and refrigeration) could be used to provide a constant and reliable grid balancing service by timing their duty cycles in response to system load.
Because it is possible to measure grid frequency from any power outlet on the grid, it is possible to design controllers for electrical appliances that detect any frequency imbalance in real time. Dynamic-demand enabled appliances would react to this same signal. When the frequency decreases they would be more likely to switch off, reducing the load on the grid and helping to restore the balance. When the frequency increases past the standard, they would be more likely to switch on, using up the excess power. Obviously, the controller must also ensure that at no point does the appliance stray out of its acceptable operating range. As line frequency is directly related to the speed of rotation of generators on the system, millions of such devices acting together would act like a huge, fast-reacting peaking power plant.
Ancillary services
The dynamic controller could also provide other ancillary services, such as aiding blackstart recovery—the ability of a power grid to be brought back to service after a power outage – if programmed with that function. Generally blackstarts are made more difficult because of the large number of reactive loads attempting to draw power simultaneously at start up when voltages are low. This causes huge overloads that trip local breakers delaying full system recovery. The dynamic controller could have these loads "wait their turn", as it were, until full power had been restored.
Another vital balancing service is ‘fast reserve’ which is the use of standby plant to replace possible lost generation (e.g. due to a failed power generator or lost power line). By shedding load quickly while the running generators spin up, then switching back in to bring the frequency back to standard, dynamic controllers could spare the high cost of fast reserve generators. Also the fast response speed of this method would avoid possible brownouts occurring.
The technology could also help facilitate greater use of generation from variable sources, like wind power. Demand-side techniques could be an efficient and cost-effective way to help integrate this resource onto the grid. In particular it would allow these sources to work in conjunction with virtual power reserves like municipal water towers to provide a reasonably predictable dispatchable capacity.
Implementation issues
Dynamic demand devices have the potential to save considerable amounts of energy by the services they provide. But before dynamic demand control can be widely incorporated regulation must be put in place to mandate installation on at least new appliances or an effective market mechanism must be created to reward installation of the technology fairly. One method contemplated is to enable the electricity meter that measures the electricity consumption also measure the grid frequency, and switch to a higher tariff if the frequency drops below a certain level. The monthly electricity bill will then say that so many hours (and so many kilowatt hours) were on the Regular tariff and a few hours on the Short Supply tariff. Those consumers without smart demand management have to pay the extra cost, but those who install smart technologies that adapt to the short supply periods will save money.
On 1 March 2011, RLtec launched its Dynamic Demand frequency response service in hot water and HVAC load devices distributed across one of the UK's largest supermarket chains, Sainsbury's. This megawatt scale virtual power plant service provides commercial frequency regulating response to National Grid in the UK. The company is now called Open Energi.
Frequency service and reserve service
The national grid in the UK already is a massive user of this technology at an industrial scale - up to 2GW of load can be lost instantaneously by frequency sensitive relays switching off steelworks etc., which is matched over a 20-minute cycle by up to 2GW of quite small emergency diesel generators. For a complete description of this complex system see for example "Emergency Diesel Standby Generator’s Potential Contribution to Dealing With Renewable Energy Sources Intermittency And Variability" - a talk by David Andrews of Wessex Water who works closely with the UK National Grid to provide this service, given at the Open University Seminar "Coping with Variability - Integrating Renewables into the Electricity System" 24 January 2006.
Up to 5GW of such diesel generation is used in France for similar purposes, but these technologies seem to be relatively unknown. There is no reason they should not be massively increased in scope to cope with even the intermittence introduced by wind power.
UK government investigation
In August 2007, the UK government published a report outlining what potential it sees for dynamic demand technology. The report stops short of recommending the government encourage its introduction. It lists a number of technical and economic barriers to its introduction and recommends these be investigated before the government encourage the use of dynamic demand. Dynamic demand is one element of a wider government investigation into technologies that can cut greenhouse gas emissions.
However, in 2009 it was announced that domestic refrigerators are now being sold into the UK incorporating a dynamic load control system.
| Technology | Electricity transmission and distribution | null |
5072278 | https://en.wikipedia.org/wiki/Molecular%20model | Molecular model | A molecular model is a physical model of an atomistic system that represents molecules and their processes. They play an important role in understanding chemistry and generating and testing hypotheses. The creation of mathematical models of molecular properties and behavior is referred to as molecular modeling, and their graphical depiction is referred to as molecular graphics.
The term, "molecular model" refer to systems that contain one or more explicit atoms (although solvent atoms may be represented implicitly) and where nuclear structure is neglected. The electronic structure is often also omitted unless it is necessary in illustrating the function of the molecule being modeled.
Molecular models may be created for several reasons – as pedagogic tools for students or those unfamiliar with atomistic structures; as objects to generate or test theories (e.g., the structure of DNA); as analogue computers (e.g., for measuring distances and angles in flexible systems); or as aesthetically pleasing objects on the boundary of art and science.
The construction of physical models is often a creative act, and many bespoke examples have been carefully created in the workshops of science departments. There is a very wide range of approaches to physical modeling, including ball-and-stick models available for purchase commercially, to molecular models created using 3D printers. The main strategy, initially in textbooks and research articles and more recently on computers. Molecular graphics has made the visualization of molecular models on computer hardware easier, more accessible, and inexpensive, although physical models are widely used to enhance the tactile and visual message being portrayed.
History
In the 1600s, Johannes Kepler speculated on the symmetry of snowflakes and the close packing of spherical objects such as fruit. The symmetrical arrangement of closely packed spheres informed theories of molecular structure in the late 1800s, and many theories of crystallography and solid state inorganic structure used collections of equal and unequal spheres to simulate packing and predict structure.
John Dalton represented compounds as aggregations of circular atoms, and although Johann Josef Loschmidt did not create physical models, his diagrams based on circles are two-dimensional analogues of later models. August Wilhelm von Hofmann is credited with the first physical molecular model around 1860. Note how the size of the carbon appears smaller than the hydrogen. The importance of stereochemistry was not then recognised and the model is essentially topological (it should be a 3-dimensional tetrahedron).
Jacobus Henricus van 't Hoff and Joseph Le Bel introduced the concept of chemistry in three dimensions of space, that is, stereochemistry. Van 't Hoff built tetrahedral molecules representing the three-dimensional properties of carbon.
Models based on spheres
Repeating units will help to show how easy it is and clear it is to represent molecules through balls that represent atoms.
The binary compounds sodium chloride (NaCl) and caesium chloride (CsCl) have cubic structures but have different space groups. This can be rationalised in terms of close packing of spheres of different sizes. For example, NaCl can be described as close-packed chloride ions (in a face-centered cubic lattice) with sodium ions in the octahedral holes. After the development of X-ray crystallography as a tool for determining crystal structures, many laboratories built models based on spheres. With the development of plastic or polystyrene balls it is now easy to create such models.
Models based on ball-and-stick
The concept of the chemical bond as a direct link between atoms can be modelled by linking balls (atoms) with sticks/rods (bonds). This has been extremely popular and is still widely used today. Initially atoms were made of spherical wooden balls with specially drilled holes for rods. Thus carbon can be represented as a sphere with four holes at the tetrahedral angles cos−1(−) ≈ 109.47°.
A problem with rigid bonds and holes is that systems with arbitrary angles could not be built. This can be overcome with flexible bonds, originally helical springs but now usually plastic. This also allows double and triple bonds to be approximated by multiple single bonds.
The model shown to the left represents a ball-and-stick model of proline. The balls have colours: black represents carbon (C); red, oxygen (O); blue, nitrogen (N); and white, hydrogen (H). Each ball is drilled with as many holes as its conventional valence (C: 4; N: 3; O: 2; H: 1) directed towards the vertices of a tetrahedron. Single bonds are represented by (fairly) rigid grey rods. Double and triple bonds use two longer flexible bonds which restrict rotation and support conventional cis/trans stereochemistry.
However, most molecules require holes at other angles and specialist companies manufacture kits and bespoke models. Besides tetrahedral, trigonal and octahedral holes, there were all-purpose balls with 24 holes. These models allowed rotation about the single rod bonds, which could be both an advantage (showing molecular flexibility) and a disadvantage (models are floppy). The approximate scale was 5 cm per ångström (0.5 m/nm or 500,000,000:1), but was not consistent over all elements.
Arnold Beevers in Edinburgh created small models using PMMA balls and stainless steel rods. By using individually drilled balls with precise bond angles and bond lengths in these models, large crystal structures to be accurately created, but with light and rigid form. Figure 4 shows a unit cell of ruby in this style.
Skeletal models
Crick and Watson's DNA model and the protein-building kits of Kendrew were among the first skeletal models. These were based on atomic components where the valences were represented by rods; the atoms were points at the intersections. Bonds were created by linking components with tubular connectors with locking screws.
André Dreiding introduced a molecular modelling kit in the late 1950s which dispensed with the connectors. A given atom would have solid and hollow valence spikes. The solid rods clicked into the tubes forming a bond, usually with free rotation. These were and are very widely used in organic chemistry departments and were made so accurately that interatomic measurements could be made by ruler.
More recently, inexpensive plastic models (such as Orbit) use a similar principle. A small plastic sphere has protuberances onto which plastic tubes can be fitted. The flexibility of the plastic means that distorted geometries can be made.
Polyhedral models
Many inorganic solids consist of atoms surrounded by a coordination sphere of electronegative atoms (e.g. PO4 tetrahedra, TiO6 octahedra). Structures can be modelled by gluing together polyhedra made of paper or plastic.
Composite models
A good example of composite models is the Nicholson approach, widely used from the late 1970s for building models of biological macromolecules. The components are primarily amino acids and nucleic acids with preformed residues representing groups of atoms. Many of these atoms are directly moulded into the template, and fit together by pushing plastic stubs into small holes. The plastic grips well and makes bonds difficult to rotate, so that arbitrary torsion angles can be set and retain their value. The conformations of the backbone and side chains are determined by pre-computing the torsion angles and then adjusting the model with a protractor.
The plastic is white and can be painted to distinguish between O and N atoms. Hydrogen atoms are normally implicit and modelled by snipping off the spokes. A model of a typical protein with approximately 300 residues could take a month to build. It was common for laboratories to build a model for each protein solved. By 2005, so many protein structures were being determined that relatively few models were made.
Computer-based models
With the development of computer-based physical modelling, it is now possible to create complete single-piece models by feeding the coordinates of a surface into the computer. Figure 6 shows models of anthrax toxin, left (at a scale of approximately 20 Å/cm or 1:5,000,000) and green fluorescent protein, right (5 cm high, at a scale of about 4 Å/cm or 1:25,000,000) from 3D Molecular Design. Models are made of plaster or starch, using a rapid prototyping process.
It has also recently become possible to create accurate molecular models inside glass blocks using a technique known as subsurface laser engraving. The image at right shows the 3D structure of an E. coli protein (DNA polymerase beta-subunit, PDB code 1MMI) etched inside a block of glass by British company Luminorum Ltd.
Computational Models
Computers can also model molecules mathematically. Programs such as Avogadro can run on typical desktops and can predict bond lengths and angles, molecular polarity and charge distribution, and even quantum mechanical properties such as absorption and emission spectra. However, these sorts of programs cannot model molecules as more atoms are added, because the number of calculations is quadratic in the number of atoms involved; if four times as many atoms are used in a molecule, the calculations with take 16 times as long. For most practical purposes, such as drug design or protein folding, the calculations of a model require supercomputing or cannot be done on classical computers at all in a reasonable amount of time. Quantum computers can model molecules with fewer calculations because the type of calculations performed in each cycle by a quantum computer are well-suited to molecular modelling.
Common colors
Some of the most common colors used in molecular models are as follows:
{| class="wikitable"
|Hydrogen
|bgcolor=white| ||white
|-
|Alkali metals
|bgcolor=darkviolet| ||violet
|-
|Alkaline earth metals
|bgcolor=darkgreen| ||dark green
|-
|Boron, most transition metals
|bgcolor=pink| ||Pink
|-
|Carbon
|bgcolor=black| ||black
|-
|Nitrogen
|bgcolor=blue| ||blue
|-
|Oxygen
|bgcolor=red| ||red
|-
|Fluorine
|bgcolor=chartreuse| ||green yellow
|-
|Chlorine
|bgcolor=limegreen| ||lime green
|-
|Bromine
|bgcolor=darkred| ||dark red
|-
|Iodine
|bgcolor=purple| ||dark violet
|-
|Noble gases
|bgcolor=cyan| ||cyan
|-
|Phosphorus
|bgcolor=darkorange| ||orange
|-
|Sulfur
|bgcolor=yellow| ||yellow
|-
|Titanium
|bgcolor=gray| ||gray
|-
|Copper
|bgcolor=ff8866| ||apricot
|-
|Mercury
|bgcolor=lightgrey| ||light grey
|}
Chronology
This table is an incomplete chronology of events where physical molecular models provided major scientific insights.
| Physical sciences | Substance | Chemistry |
2003525 | https://en.wikipedia.org/wiki/Painite | Painite | Painite is a very rare borate mineral. It was first found in Myanmar by British mineralogist and gem dealer Arthur C.D. Pain who misidentified it as ruby, until it was discovered as a new gemstone in the 1950s. When it was confirmed as a new mineral species, the mineral was named after him.
The chemical makeup of painite contains calcium, zirconium, boron, aluminium, and oxygen (CaZrAl9O15(BO3)). The mineral also contains trace amounts of chromium and vanadium, which are responsible for Painite's typically orange-red to brownish-red color, similar to topaz. The mineral's rarity is due to zirconium and boron rarely interacting with each other in nature. The crystals are naturally hexagonal, but may also be euhedral or orthorhombic. They also may have no crystalline structure, but usually are accompanied by a crystalline structure. Until late 2004, only two had been cut into faceted gemstones.
Discovery and occurrence
Extensive exploration in the area surrounding Mogok, which comprises a large part of the extremely small region the mineral is known to exist in, has identified several new painite occurrences that have been vigorously explored resulting in several thousand new available painite specimens.
| Physical sciences | Minerals | Earth science |
3744388 | https://en.wikipedia.org/wiki/Ohmic%20contact | Ohmic contact | An ohmic contact is a non-rectifying electrical junction: a junction between two conductors that has a linear current–voltage (I–V) curve as with Ohm's law. Low-resistance ohmic contacts are used to allow charge to flow easily in both directions between the two conductors, without blocking due to rectification or excess power dissipation due to voltage thresholds.
By contrast, a junction or contact that does not demonstrate a linear I–V curve is called non-ohmic. Non-ohmic contacts come in a number of forms, such as p–n junction, Schottky barrier, rectifying heterojunction, or breakdown junction.
Generally the term "ohmic contact" implicitly refers to an ohmic contact of a metal to a semiconductor, where achieving ohmic contact resistance is possible but requires careful technique. Metal–metal ohmic contacts are relatively simpler to make, by ensuring direct contact between the metals without intervening layers of insulating contamination, excessive roughness or oxidation; various techniques are used to create ohmic metal–metal junctions (soldering, welding, crimping, deposition, electroplating, etc.). This article focuses on metal–semiconductor ohmic contacts.
Stable contacts at semiconductor interfaces, with low contact resistance and linear I–V behavior, are critical for the performance and reliability of semiconductor devices, and their preparation and characterization are major efforts in circuit fabrication. Poorly prepared junctions to semiconductors can easily show rectifying behaviour by causing depletion of the semiconductor near the junction, rendering the device useless by blocking the flow of charge between those devices and the external circuitry. Ohmic contacts to semiconductors are typically constructed by depositing thin metal films of a carefully chosen composition, possibly followed by annealing to alter the semiconductor–metal bond.
Physics of formation of metal–semiconductor ohmic contacts
Both ohmic contacts and Schottky barriers are dependent on the Schottky barrier height, which sets the threshold for the excess energy an electron requires to pass from the semiconductor to the metal. For the junction to admit electrons easily in both directions (ohmic contact), the barrier height must be small in at least some parts of the junction surface. To form an excellent ohmic contact (low resistance), the barrier height should be small everywhere and furthermore the interface should not reflect electrons.
The Schottky barrier height between a metal and semiconductor is naively predicted by the Schottky–Mott rule to be proportional to the difference of the metal-vacuum work function and the semiconductor-vacuum electron affinity. In practice, most metal–semiconductor interfaces do not follow this rule to the predicted degree. Instead, the chemical termination of the semiconductor crystal against a metal creates electron states within its band gap. The nature of these metal-induced gap states and their occupation by electrons tends to pin the center of the band gap to the Fermi level, an effect known as Fermi level pinning. Thus, the heights of the Schottky barriers in metal–semiconductor contacts often show little dependence on the value of the semiconductor or metal work functions, in stark contrast to the Schottky–Mott rule. Different semiconductors exhibit this Fermi level pinning to different degrees, but a technological consequence is that high quality (low resistance) ohmic contacts are usually difficult to form in important semiconductors such as silicon and gallium arsenide.
The Schottky–Mott rule is not entirely incorrect since, in practice, metals with high work functions form the best contacts to p-type semiconductors, while those with low work functions form the best contacts to n-type semiconductors. Unfortunately experiments have shown that the predictive power of the model doesn't extend much beyond this statement. Under realistic conditions, contact metals may react with semiconductor surfaces to form a compound with new electronic properties. A contamination layer at the interface may effectively widen the barrier. The surface of the semiconductor may reconstruct leading to a new electronic state. The dependence of contact resistance on the details of the interfacial chemistry is what makes the reproducible fabrication of ohmic contacts such a manufacturing challenge.
Preparation and characterization of ohmic contacts
The fabrication of the ohmic contacts is a much-studied part of materials engineering that nonetheless remains something of an art. The reproducible, reliable fabrication of contacts relies on extreme cleanliness of the semiconductor surface. Since a native oxide rapidly forms on the surface of silicon, for example, the performance of a contact can depend sensitively on the details of preparation.
Often the contact region is heavily doped to ensure the type of contact wanted. As a rule, ohmic contacts on semiconductors form more easily when the semiconductor is highly doped near the junction; a high doping narrows the depletion region at the interface and allow electrons to flow in both directions easily at any bias by tunneling through the barrier.
The fundamental steps in contact fabrication are semiconductor surface cleaning, contact metal deposition, patterning and annealing. Surface cleaning may be performed by sputter-etching, chemical etching, reactive gas etching or ion milling. For example, the native oxide of silicon may be removed with a hydrofluoric acid dip, while GaAs is more typically cleaned by a bromine-methanol dip. After cleaning, metals are deposited via sputter deposition, evaporation or chemical vapor deposition (CVD). Sputtering is a faster and more convenient method of metal deposition than evaporation but the ion bombardment from the plasma may induce surface states or even invert the charge carrier type at the surface. For this reason the gentler but still rapid CVD may be used. Post-deposition annealing of contacts is useful for relieving stress as well as for inducing any desirable reactions between the metal and the semiconductor.
Because deposited metals can themselves react in ambient conditions, to the detriment of the contacts' electrical properties, it is common to form ohmic contacts with layered structures, with the bottom layer, in contact with the semiconductor, chosen for its ability to induce ohmic behaviour. A diffusion barrier-layer may be used to prevent the layers from mixing during any annealing process.
The measurement of contact resistance is most simply performed using a four-point probe although for more accurate determination, use of the transmission line method is typical.
Technologically important kinds of contacts
Aluminum was originally the most important contact metal for silicon which was used with either the n-type or p-type semiconductor. As with other reactive metals, Al contributes to contact formation by consuming oxygen from native silicon-dioxide residue. Pure aluminum did react with the silicon, so it was replaced by silicon-doped aluminum and eventually by silicides less prone to diffuse during subsequent high-temperature processing.
Modern ohmic contacts to silicon such as titanium-tungsten disilicide are usually silicides made by CVD. Contacts are often made by depositing the transition metal and forming the silicide by annealing with the result that the silicide may be non-stoichiometric. Silicide contacts can also be deposited by direct sputtering of the compound or by ion implantation of the transition metal followed by annealing.
Formation of contacts to compound semiconductors is considerably more difficult than with silicon. For example, GaAs surfaces tend to lose arsenic and the trend towards As loss can be considerably exacerbated by the deposition of metal. In addition, the volatility of As limits the amount of post-deposition annealing that GaAs devices will tolerate. One solution for GaAs and other compound semiconductors is to deposit a low-bandgap alloy contact layer as opposed to a heavily doped layer. For example, GaAs itself has a smaller bandgap than AlGaAs and so a layer of GaAs near its surface can promote ohmic behavior. In general the technology of ohmic contacts for III-V and II-VI semiconductors is much less developed than for Si.
Transparent or semi-transparent contacts are necessary for active matrix LCD displays, optoelectronic devices such as laser diodes and photovoltaics. The most popular choice is indium tin oxide, a metal that is formed by reactive sputtering of an In-Sn target in an oxide atmosphere.
Significance
The RC time constant associated with contact resistance can limit the frequency response of devices. The charging and discharging of the leads resistance is a major cause of power dissipation in high-clock-rate digital electronics. Contact resistance causes power dissipation by Joule heating in low-frequency and analog circuits (for example, solar cells) made from less common semiconductors. The establishment of a contact fabrication methodology is a critical part of the technological development of any new semiconductor. Electromigration and delamination at contacts are also a limitation on the lifetime of electronic devices.
| Physical sciences | Electrical circuits | Physics |
25652303 | https://en.wikipedia.org/wiki/Computer%20architecture | Computer architecture | In computer science and computer engineering, computer architecture is a description of the structure of a computer system made from component parts. It can sometimes be a high-level description that ignores details of the implementation. At a more detailed level, the description may include the instruction set architecture design, microarchitecture design, logic design, and implementation.
History
The first documented computer architecture was in the correspondence between Charles Babbage and Ada Lovelace, describing the analytical engine. While building the computer Z1 in 1936, Konrad Zuse described in two patent applications for his future projects that machine instructions could be stored in the same storage used for data, i.e., the stored-program concept. Two other early and important examples are:
John von Neumann's 1945 paper, First Draft of a Report on the EDVAC, which described an organization of logical elements; and
Alan Turing's more detailed Proposed Electronic Calculator for the Automatic Computing Engine, also 1945 and which cited John von Neumann's paper.
The term "architecture" in computer literature can be traced to the work of Lyle R. Johnson and Frederick P. Brooks, Jr., members of the Machine Organization department in IBM's main research center in 1959. Johnson had the opportunity to write a proprietary research communication about the Stretch, an IBM-developed supercomputer for Los Alamos National Laboratory (at the time known as Los Alamos Scientific Laboratory). To describe the level of detail for discussing the luxuriously embellished computer, he noted that his description of formats, instruction types, hardware parameters, and speed enhancements were at the level of "system architecture", a term that seemed more useful than "machine organization".
Subsequently, Brooks, a Stretch designer, opened Chapter 2 of a book called Planning a Computer System: Project Stretch by stating, "Computer architecture, like other architecture, is the art of determining the needs of the user of a structure and then designing to meet those needs as effectively as possible within economic and technological constraints."
Brooks went on to help develop the IBM System/360 line of computers, in which "architecture" became a noun defining "what the user needs to know". The System/360 line was succeeded by several compatible lines of computers, including the current IBM Z line. Later, computer users came to use the term in many less explicit ways.
The earliest computer architectures were designed on paper and then directly built into the final hardware form.
Later, computer architecture prototypes were physically built in the form of a transistor–transistor logic (TTL) computer—such as the prototypes of the 6800 and the PA-RISC—tested, and tweaked, before committing to the final hardware form.
As of the 1990s, new computer architectures are typically "built", tested, and tweaked—inside some other computer architecture in a computer architecture simulator; or inside a FPGA as a soft microprocessor; or both—before committing to the final hardware form.
Subcategories
The discipline of computer architecture has three main subcategories:
Instruction set architecture (ISA): defines the machine code that a processor reads and acts upon as well as the word size, memory address modes, processor registers, and data type.
Microarchitecture: also known as "computer organization", this describes how a particular processor will implement the ISA. The size of a computer's CPU cache for instance, is an issue that generally has nothing to do with the ISA.
Systems design: includes all of the other hardware components within a computing system, such as data processing other than the CPU (e.g., direct memory access), virtualization, and multiprocessing.
There are other technologies in computer architecture. The following technologies are used in bigger companies like Intel, and were estimated in 2002 to count for 1% of all of computer architecture:
Macroarchitecture: architectural layers more abstract than microarchitecture
Assembly instruction set architecture: A smart assembler may convert an abstract assembly language common to a group of machines into slightly different machine language for different implementations.
Programmer-visible macroarchitecture: higher-level language tools such as compilers may define a consistent interface or contract to programmers using them, abstracting differences between underlying ISAs and microarchitectures. For example, the C, C++, or Java standards define different programmer-visible macroarchitectures.
Microcode: microcode is software that translates instructions to run on a chip. It acts like a wrapper around the hardware, presenting a preferred version of the hardware's instruction set interface. This instruction translation facility gives chip designers flexible options: E.g. 1. A new improved version of the chip can use microcode to present the exact same instruction set as the old chip version, so all software targeting that instruction set will run on the new chip without needing changes. E.g. 2. Microcode can present a variety of instruction sets for the same underlying chip, allowing it to run a wider variety of software.
Pin architecture: The hardware functions that a microprocessor should provide to a hardware platform, e.g., the x86 pins A20M, FERR/IGNNE or FLUSH. Also, messages that the processor should emit so that external caches can be invalidated (emptied). Pin architecture functions are more flexible than ISA functions because external hardware can adapt to new encodings, or change from a pin to a message. The term "architecture" fits, because the functions must be provided for compatible systems, even if the detailed method changes.
Roles
Definition
Computer architecture is concerned with balancing the performance, efficiency, cost, and reliability of a computer system. The case of instruction set architecture can be used to illustrate the balance of these competing factors. More complex instruction sets enable programmers to write more space efficient programs, since a single instruction can encode some higher-level abstraction (such as the x86 Loop instruction). However, longer and more complex instructions take longer for the processor to decode and can be more costly to implement effectively. The increased complexity from a large instruction set also creates more room for unreliability when instructions interact in unexpected ways.
The implementation involves integrated circuit design, packaging, power, and cooling. Optimization of the design requires familiarity with topics from compilers and operating systems to logic design and packaging.
Instruction set architecture
An instruction set architecture (ISA) is the interface between the computer's software and hardware and also can be viewed as the programmer's view of the machine. Computers do not understand high-level programming languages such as Java, C++, or most programming languages used. A processor only understands instructions encoded in some numerical fashion, usually as binary numbers. Software tools, such as compilers, translate those high level languages into instructions that the processor can understand.
Besides instructions, the ISA defines items in the computer that are available to a program—e.g., data types, registers, addressing modes, and memory. Instructions locate these available items with register indexes (or names) and memory addressing modes.
The ISA of a computer is usually described in a small instruction manual, which describes how the instructions are encoded. Also, it may define short (vaguely) mnemonic names for the instructions. The names can be recognized by a software development tool called an assembler. An assembler is a computer program that translates a human-readable form of the ISA into a computer-readable form. Disassemblers are also widely available, usually in debuggers and software programs to isolate and correct malfunctions in binary computer programs.
ISAs vary in quality and completeness. A good ISA compromises between programmer convenience (how easy the code is to understand), size of the code (how much code is required to do a specific action), cost of the computer to interpret the instructions (more complexity means more hardware needed to decode and execute the instructions), and speed of the computer (with more complex decoding hardware comes longer decode time). Memory organization defines how instructions interact with the memory, and how memory interacts with itself.
During design emulation, emulators can run programs written in a proposed instruction set. Modern emulators can measure size, cost, and speed to determine whether a particular ISA is meeting its goals.
Computer organization
Computer organization helps optimize performance-based products. For example, software engineers need to know the processing power of processors. They may need to optimize software in order to gain the most performance for the lowest price. This can require quite a detailed analysis of the computer's organization. For example, in an SD card, the designers might need to arrange the card so that the most data can be processed in the fastest possible way.
Computer organization also helps plan the selection of a processor for a particular project. Multimedia projects may need very rapid data access, while virtual machines may need fast interrupts. Sometimes certain tasks need additional components as well. For example, a computer capable of running a virtual machine needs virtual memory hardware so that the memory of different virtual computers can be kept separated. Computer organization and features also affect power consumption and processor cost.
Implementation
Once an instruction set and microarchitecture have been designed, a practical machine must be developed. This design process is called the implementation. Implementation is usually not considered architectural design, but rather hardware design engineering. Implementation can be further broken down into several steps:
Logic implementation designs the circuits required at a logic-gate level.
Circuit implementation does transistor-level designs of basic elements (e.g., gates, multiplexers, latches) as well as of some larger blocks (ALUs, caches etc.) that may be implemented at the logic-gate level, or even at the physical level if the design calls for it.
Physical implementation draws physical circuits. The different circuit components are placed in a chip floor plan or on a board and the wires connecting them are created.
Design validation tests the computer as a whole to see if it works in all situations and all timings. Once the design validation process starts, the design at the logic level are tested using logic emulators. However, this is usually too slow to run a realistic test. So, after making corrections based on the first test, prototypes are constructed using Field-Programmable Gate-Arrays (FPGAs). Most hobby projects stop at this stage. The final step is to test prototype integrated circuits, which may require several redesigns.
For CPUs, the entire implementation process is organized differently and is often referred to as CPU design.
Design goals
The exact form of a computer system depends on the constraints and goals. Computer architectures usually trade off standards, power versus performance, cost, memory capacity, latency (latency is the amount of time that it takes for information from one node to travel to the source) and throughput. Sometimes other considerations, such as features, size, weight, reliability, and expandability are also factors.
The most common scheme does an in-depth power analysis and figures out how to keep power consumption low while maintaining adequate performance.
Performance
Modern computer performance is often described in instructions per cycle (IPC), which measures the efficiency of the architecture at any clock frequency; a faster IPC rate means the computer is faster. Older computers had IPC counts as low as 0.1 while modern processors easily reach nearly 1. Superscalar processors may reach three to five IPC by executing several instructions per clock cycle.
Counting machine-language instructions would be misleading because they can do varying amounts of work in different ISAs. The "instruction" in the standard measurements is not a count of the ISA's machine-language instructions, but a unit of measurement, usually based on the speed of the VAX computer architecture.
Many people used to measure a computer's speed by the clock rate (usually in MHz or GHz). This refers to the cycles per second of the main clock of the CPU. However, this metric is somewhat misleading, as a machine with a higher clock rate may not necessarily have greater performance. As a result, manufacturers have moved away from clock speed as a measure of performance.
Other factors influence speed, such as the mix of functional units, bus speeds, available memory, and the type and order of instructions in the programs.
There are two main types of speed: latency and throughput. Latency is the time between the start of a process and its completion. Throughput is the amount of work done per unit time. Interrupt latency is the guaranteed maximum response time of the system to an electronic event (like when the disk drive finishes moving some data).
Performance is affected by a very wide range of design choices — for example, pipelining a processor usually makes latency worse, but makes throughput better. Computers that control machinery usually need low interrupt latencies. These computers operate in a real-time environment and fail if an operation is not completed in a specified amount of time. For example, computer-controlled anti-lock brakes must begin braking within a predictable and limited time period after the brake pedal is sensed or else failure of the brake will occur.
Benchmarking takes all these factors into account by measuring the time a computer takes to run through a series of test programs. Although benchmarking shows strengths, it should not be how you choose a computer. Often the measured machines split on different measures. For example, one system might handle scientific applications quickly, while another might render video games more smoothly. Furthermore, designers may target and add special features to their products, through hardware or software, that permit a specific benchmark to execute quickly but do not offer similar advantages to general tasks.
Power efficiency
Power efficiency is another important measurement in modern computers. Higher power efficiency can often be traded for lower speed or higher cost. The typical measurement when referring to power consumption in computer architecture is MIPS/W (millions of instructions per second per watt).
Modern circuits have less power required per transistor as the number of transistors per chip grows. This is because each transistor that is put in a new chip requires its own power supply and requires new pathways to be built to power it. However, the number of transistors per chip is starting to increase at a slower rate. Therefore, power efficiency is starting to become as important, if not more important than fitting more and more transistors into a single chip. Recent processor designs have shown this emphasis as they put more focus on power efficiency rather than cramming as many transistors into a single chip as possible. In the world of embedded computers, power efficiency has long been an important goal next to throughput and latency.
Shifts in market demand
Increases in clock frequency have grown more slowly over the past few years, compared to power reduction improvements. This has been driven by the end of Moore's Law and demand for longer battery life and reductions in size for mobile technology. This change in focus from higher clock rates to power consumption and miniaturization can be shown by the significant reductions in power consumption, as much as 50%, that were reported by Intel in their release of the Haswell microarchitecture; where they dropped their power consumption benchmark from 30–40 watts down to 10–20 watts. Comparing this to the processing speed increase of 3 GHz to 4 GHz (2002 to 2006), it can be seen that the focus in research and development is shifting away from clock frequency and moving towards consuming less power and taking up less space.
| Technology | Computer science | null |
25658445 | https://en.wikipedia.org/wiki/Mass%20concentration%20%28chemistry%29 | Mass concentration (chemistry) | In chemistry, the mass concentration (or ) is defined as the mass of a constituent divided by the volume of the mixture .
For a pure chemical the mass concentration equals its density (mass divided by volume); thus the mass concentration of a component in a mixture can be called the density of a component in a mixture. This explains the usage of (the lower case Greek letter rho), the symbol most often used for density.
Definition and properties
The volume in the definition refers to the volume of the solution, not the volume of the solvent. One litre of a solution usually contains either slightly more or slightly less than 1 litre of solvent because the process of dissolution causes volume of liquid to increase or decrease. Sometimes the mass concentration is called titre.
Notation
The notation common with mass density underlines the connection between the two quantities (the mass concentration being the mass density of a component in the solution), but it can be a source of confusion especially when they appear in the same formula undifferentiated by an additional symbol (like a star superscript, a bolded symbol or varrho).
Dependence on volume
Mass concentration depends on the variation of the volume of the solution due mainly to thermal expansion. On small intervals of temperature the dependence is :
where is the mass concentration at a reference temperature, is the thermal expansion coefficient of the mixture.
Sum of mass concentrations - normalizing relation
The sum of the mass concentrations of all components (including the solvent) gives the density of the solution:
Thus, for pure component the mass concentration equals the density of the pure component.
Units
The SI-unit for mass concentration is kg/m3 (kilogram/cubic metre). This is the same as mg/mL and g/L. Another commonly used unit is g/(100 mL), which is identical to g/dL (gram/decilitre).
Usage in biology
In biology and medicine, the "%" symbol is widely used in a misnomer sense to denote mass concentration, also called "mass/volume percentage". A solution with 1 g of solute dissolved in a final volume of 100 mL of solution would be labeled as "1%" or "1% m/v" (mass/volume). The common names of intravenous sugar solutions, such as D5W and D50W, reflect this convention. The notation is mathematically flawed because the unit "%" can only be used for dimensionless quantities. "Percent solution" or "percentage solution" are thus terms best reserved for "mass percent solutions" (m/m = m% = mass solute/mass total solution after mixing), or "volume percent solutions" (v/v = v% = volume solute per volume of total solution after mixing). The very ambiguous terms "percent solution" and "percentage solutions" with no other qualifiers, continue to occasionally be encountered.
This common usage of % to mean m/v in biology is because of many biological solutions being dilute and water-based, an aqueous solution. Liquid water has a density of approximately 1 g/cm3 (1 g/mL). Thus 100 mL of water is equal to approximately 100 g. Therefore, a solution with 1 g of solute dissolved in final volume of 100 mL aqueous solution may also be considered 1% m/m (1 g solute in 99 g water). This approximation breaks down as the solute concentration is increased (for example, in water–NaCl mixtures). High solute concentrations are often not physiologically relevant, but are occasionally encountered in pharmacology, where the mass per volume notation is still sometimes encountered. An extreme example is saturated solution of potassium iodide (SSKI) which attains 100 "%" m/v potassium iodide mass concentration (1 gram KI per 1 mL solution) only because the solubility of the dense salt KI is extremely high in water, and the resulting solution is very dense (1.72 times as dense as water).
Although there are examples to the contrary, it should be stressed that the commonly used "units" of % w/v are grams per millilitre (g/mL). 1% m/v solutions are sometimes thought of as being gram/100 mL but this detracts from the fact that % m/v is g/mL; 1g of water has a volume of approximately 1 mL (at standard temperature and pressure) and the mass concentration is said to be 100%. To make 10 mL of an aqueous 1% cholate solution, 0.1 grams of cholate are dissolved in 10mL of water. Volumetric flasks are the most appropriate piece of glassware for this procedure as deviations from ideal solution behavior can occur with high solute concentrations.
In solutions, mass concentration is commonly encountered as the ratio of mass/[volume solution], or m/v. In water solutions containing relatively small quantities of dissolved solute (as in biology), such figures may be "percentivized" by multiplying by 100 a ratio of grams solute per mL solution. The result is given as "mass/volume percentage". Such a convention expresses mass concentration of 1 gram of solute in 100 mL of solution, as "1 m/v %".
Related quantities
Density of pure component
The relation between mass concentration and density of a pure component (mass concentration of single component mixtures) is:
where is the density of the pure component, the volume of the pure component before mixing.
Specific volume (or mass-specific volume)
Specific volume is the inverse of mass concentration only in the case of pure substances, for which mass concentration is the same as the density of the pure-substance:
Molar concentration
The conversion to molar concentration is given by:
where is the molar mass of constituent .
Mass fraction
The conversion to mass fraction is given by:
Mole fraction
The conversion to mole fraction is given by:
where is the average molar mass of the mixture.
Molality
For binary mixtures, the conversion to molality is given by:
Spatial variation and gradient
The values of (mass and molar) concentration different in space triggers the phenomenon of diffusion.
| Physical sciences | Mixture | Chemistry |
31352816 | https://en.wikipedia.org/wiki/Cobble%20%28geology%29 | Cobble (geology) | A cobble (sometimes a cobblestone) is a clast of rock defined on the Udden–Wentworth scale as having a particle size of , larger than a pebble and smaller than a boulder. Other scales define a cobble's size differently. A rock made predominantly of cobbles is termed a conglomerate. Cobblestone is a building material based on cobbles.
Etymology
Cobbles, also called cobblestones, derive their name from the word cob, meaning a rounded lump. The term is further related to the German , meaning head. Chester Wentworth referred to cobbles as cobble bowlders in his 1922 paper that would become the basis for the Udden–Wentworth scale.
Classifications
Within the widely used Krumbein phi scale of grain sizes, cobbles are defined as clasts of rock ranging from −6 to −8 φ. This classification corresponds with the Udden–Wentworth size scale which defines cobbles as clasts with diameters from . On this scale, cobbles are larger than pebbles which measure in diameter and smaller than boulders, whose diameters range from . On the Udden–Wentworth scale, an unlithified fraction of cobbles is classified as gravel while a lithified sample primarily composed of cobbles is a conglomerate. The Committee on Sedimentation of the US National Research Council has recommended that in situ cobbles be identified by their process of origination, if possible (e.g., cobbles by disintegration, by exfoliation, etc.).
In the late 1800s and early to mid-1900s, prior to the Udden–Wentworth scale's widespread adoption, size classifications tended to group all particles larger than together as gravel or stones. Other scales have defined the size of a cobble slightly differently than the Udden–Wentworth; the British Standards Institution denotes a cobble as any clast ranging in diameter from while the United States Department of Agriculture's definition suggests a range of and the ISO standard 14688 names cobbles as ranging from in diameter.
Various attempts have been made to refine the Udden–Wentworth scale, including its definition of cobbles. In 1968, D. J. Doeglas proposed subdividing the cobble designation into two fractions, small cobbles (for particles with diameters from ) and large cobbles (for particles with diameters from ). A 1999 paper by Terence C. Blair and John G. McPherson argued that the Udden–Wentworth and Krumbein scales betrayed a historical emphasis on the study of sand grains while ignoring larger gravel grains. They proposed defining fine cobbles as those with diameters from (−6 to −7 φ) and coarse cobbles as those with diameters from (−7 to −8 φ). In 2012, Simon J. Blott and Kenneth Pye suggested that the cobble designation be eliminated altogether, replaced by very small boulder and small boulder designations equivalent in size to Blair and McPherson's fine and coarse cobbles, respectively.
Settings
When occurring in streams, cobbles are likely to be found in mountain valley streambeds that are moderately steep. Cobbles are also transported by glaciers and deposited as with other grades of sediment as till. If the till is water-laid, finer particles like sand and pebbles may be entirely washed away, leaving a deposit of only boulders and cobbles. The term shingle beach refers to a beach covered with small- to medium-sized cobbles or pebbles (as opposed to fine sand). Glacially transported cobbles tend to share several identifying features including a tabular shape and downward diagonal striations on lateral facets.
Cobble conglomerates may be alluvial in origin or the product of "stone avalanches", a type of debris flow resulting from unconsolidated cobbles and gravel. In such stone avalanches, well-rounded cobbles may travel the farthest on account of their low rolling friction. When the product of alluvial processes, the cobble conglomerate's matrix consists of gravel and coarse sand. In contrast, the matrices of flow-deposited conglomerates are primarily mud.
| Physical sciences | Sedimentology | Earth science |
891641 | https://en.wikipedia.org/wiki/Mexico%20City%20Metro | Mexico City Metro | The Mexico City Metro () is a rapid transit system that serves the metropolitan area of Mexico City, including some municipalities in the State of Mexico. Operated by the Sistema de Transporte Colectivo (STC), it is the second largest metro system in North America after the New York City Subway.
The inaugural STC Metro line was long, serving 16 stations, and opened to the public on 4 September 1969. The system has expanded since then in a series of fits and starts. , the system has 12 lines, serving 195 stations, and of route. Ten of the lines are rubber-tired. Instead of traditional steel wheels, they use pneumatic traction, which is quieter and rides smoother in Mexico City's unstable soils. The system survived the 1985 Mexico City earthquake.
Of the STC Metro's 195 stations, 44 serve two or more lines (correspondencias or transfer stations). Many stations are named for historical figures, places, or events in Mexican history. It has 115 underground stations (the deepest of which are below street level); 54 surface stations and 26 elevated stations. All lines operate from 5 a.m. to midnight. At the end of 2007, the Federal District government announced the construction of the most recent STC Metro line, Line 12, which was built to run approximately towards the southeastern part of the city, connecting with Lines 7, 3, 2 and 8. This line opened on 30 October 2012.
History
Concept of the Metro and early plans
By the second half of the twentieth century, Mexico City had serious public transport issues, with congested main roads and highways, especially in the downtown zone, where 40 percent of the daily trips in the city were concentrated. 65 of the 91 lines of bus and electric transport served this area. With four thousand units in addition to 150,000 personal automobile peak hours, the average speed was less than walking pace.
The principal promoter of the construction of the Mexico City Metro was engineer Bernardo Quintana, who was in charge of the construction company Ingenieros Civiles y Asociados (Civil Engineers and Associates). He carried out a series of studies that resulted in a draft plan which would ultimately lead to the construction of the Mexico City Metro. This plan was shown to different authorities of Mexico City but it was not made official until 29 April 1967, when the Government Gazette ("Diario Oficial de la Federación") published the presidential decree that created a public decentralized organism, the Sistema de Transporte Colectivo, with the proposal to build, operate and run an underground rapid transit network as part of Mexico City's public transport system.
The Mexico City Metro benefited from a great amount of technical assistance made available by France. RATP's engineering branch SOFRETU played a major role in its initial planning and the design of the first lines, hence the choice of tyre/rail technology.
On 19 June 1967, at the crossroads of Chapultepec Avenue with Avenida Bucareli, the inauguration ceremony for the Mexico City Metro took place. Two years later, on 4 September 1969, an orange train made the inaugural trip between Zaragoza and Insurgentes stations, thus beginning daily operation up to today.
First stage (1967–1972)
The first stage of construction comprised the construction, done by Grupo ICA, and inauguration of lines 1, 2 and 3. This stage involved engineers, geologists, mechanics, civil engineers, chemists, hydraulic and sanitation workers, electricians, archaeologists, and biologists; specialists in ventilation, statistics, computation, and in traffic and transit; accountants, economists, lawyers, workers and laborers. Between 1,200 and 4,000 specialists and 48,000 workers participated, building at least of track per month, the fastest rate of construction ever for a subway.
During this stage of construction workers uncovered two archaeological ruins, one Aztec idol, and the bones of a mammoth (on display at Talismán station).
By the end of the first stage, namely on 10 June 1972, the STC Metro had 48 stations and a total length of : Line 1 ran from Observatorio to Zaragoza, Line 2 from Tacuba southwest to Tasqueña and line 3 from Tlatelolco to Hospital General in the south, providing quick access to the General Hospital of Mexico.
Second stage (1977–1982)
No further progress was reached during President Luis Echeverría's government, but during José López Portillo's administration, a second stage began. The Comisión Ejecutiva del Metro (Executive Technical Commission of Mexico City Metro) was created in order to be in charge of expanding the STC Metro within the metropolitan area of Mexico City.
Works began with the expansion of Line 3 towards the north from Tlatelolco to La Raza in 1978 and to the current terminal Indios Verdes in 1979, and towards the south from Hospital General to Centro Médico in 1980 and to Zapata months later. Construction of lines 4 and 5 was begun and completed on 26 May – 30 August 1982, respectively; the former from Martín Carrera to Santa Anita and the latter from Politécnico to Pantitlán. Line 4 was the first STC Metro line built as an elevated track, owing to the lower density of big buildings.
Third stage (1983–1985), and the 1985 earthquake
This construction stage took place from the beginning of 1983 through the end of 1985. Lines 1, 2 and 3 were expanded to their current lengths, and new lines 6 and 7 were built. The length of the network was increased by and the number of stations to 105.
Line 3's route was expanded from Zapata station to Universidad station on 30 August 1983. Line 1 was expanded from Zaragoza to the current terminal Pantitlán, and line 2 from Tacuba to the current terminal Cuatro Caminos. These last two were both inaugurated on 22 August 1984.
Line 6's route first ran from El Rosario to Instituto del Petróleo; Line 7 was opened from Tacuba to Barranca del Muerto and runs along the foot of the Sierra de las Cruces mountain range that surrounds the Valley of Mexico at its west side, outside of the ancient lake zone. This made it possible for Line 7 to be built as a deep-bore tunnel.
On the morning of 19 September 1985, a magnitude 8.0 earthquake struck Mexico City. Many buildings as well as streets were left with major damage making transportation on the ground difficult, but the STC Metro was not damaged because a rectangular structure had been used instead of arches, making it resistant to earthquakes, thus proving to be a safe means of transportation in a time of crisis.
On the day of the quake, the Metro stopped service and completely shut down for fear of electrocution. This caused people to get out of the tunnels from wherever they were and onto the street to try to get where they were going. At the time, the Metro had 101 stations, with 32 closed to the public in the weeks after the event. On Line 1, there was no service in stations Merced, Pino Suárez, Isabel la Católica, Salto del Agua, Balderas or Cuauhtémoc. On Line 2, there was no service between stations Bellas Artes and Tasqueña. On Line 3 only Juárez and Balderas were closed. Line 4 continued to operate normally. All of the closed stations were in the historic center area, with the exception of the stations of Line 2 south of Pino Suárez. These stations were located above the ground. The reason these stations were closed was not due to damage to the Metro proper, but rather because of surface rescue work and clearing of debris.
Fourth stage (1985–1987)
Fourth stage saw the completion of Line 6 from Instituto del Petróleo to its eastern terminal Martín Carrera and Line 7 to the north from Tacuba to El Rosario. Line 9 was the only new line built during this stage. It originally ran from Pantitlán to Centro Médico, and its expansion to Tacubaya was completed on 29 August 1988. For Line 9, a circular deep-bore tunnel and an elevated track were used.
Fifth stage (1988–1994)
For the first time, a service line of the Mexico City Metro ran into the State of Mexico: planned as one of more líneas alimentadoras (feeder lines to be named by letters, instead of numbers), line A was fully operational by its first inauguration on 12 August 1991. It runs from Pantitlán to La Paz, located in the municipality of the same name. This line was built almost entirely above ground, and to reduce the cost of maintenance, steel railway tracks and overhead lines were used instead of pneumatic traction, promoting the name metro férreo (steel-rail metro) as opposed to the previous eight lines that used pneumatic traction.
The draft for Line 8 planned a correspondencia (transfer station) in Zócalo, namely the exact center of the city, but it was canceled due to possible damage to the colonial buildings and the Aztec ruins, so it was replanned and now it runs from Garibaldi, which is still downtown, to Constitución de 1917 in the southeast of the city. The construction of line 8 began in 1988 and was completed in 1994.
With this, the length of the network increased , adding two lines and 29 more stations, giving the metro network at that point a total of , 154 stations and 10 lines.
Sixth stage (1994–2000)
Assessment for line B began in late 1993. Line B was intended as a second línea alimentadora for northeastern municipalities in the State of Mexico, but, unlike line A, it used pneumatic traction. Construction of the subterranean track between Buenavista (named after the old Buenavista train station) and Garibaldi began in October 1994.
Line B was opened to the public in two stages: from Buenavista to Villa de Aragón on 15 December 1999, and from Villa de Aragón to Ciudad Azteca on 30 November 2000.
Seventh stage (2008–2014)
Plans for a new STC Metro line started in 2008, although previous surveys and assessments were made as early as 2000. Line 12's first service stage was planned for completion in late 2009 with the creation of track connecting Axomulco, a planned new transfer station for Line 8 (between Escuadrón 201 and Atlalilco) to Tláhuac. The second stage, connecting Mixcoac to Tláhuac, was to be completed in 2010.
Construction of Line 12 started in 2008, assuring it would be opened by 2011. Nevertheless, completion was delayed until 2012. Free test rides were offered to the public in some stations, and the line was fully operational on 30 October 2012. With minor changes, Line 12 runs from Mixcoac to Tláhuac, serving southern Mexico City for the first time. At long, it is the longest line in the system.
Line 12 differs from previous lines in several aspects: no hawkers are allowed, either inside the train or inside the stations; it is the first numbered line to use steel railway tracks; one must have a Tarjeta DF smart card to access any station since Metro tickets are no longer accepted.
In the book Los hombres del Metro, the original planning of Line 12 is described; although it was to begin at Mixcoac as it does today, Atlalilco and Constitución de 1917 stations of Line 8 were to be part of Line 12. The same map shows that Line 8 would have reached the Villa Coapa area and that it would not have had a terminal at Garibaldi, but at Indios Verdes, linking with Line 3. In addition, the book shows that Line 7 would have terminated at San Jerónimo. None of these plans have been confirmed by the Mexico City government.
In 2015, mayor Miguel Ángel Mancera announced the construction of two more stations and a terminal for Line 12: Valentín Campa, Álvaro Obregón and Observatorio, both west of Mixcoac. With this, Line 12 is to be connected to Line 1, providing new metro access to the Observatorio zone, which will become the terminal for the intercity train between Mexico City and Toluca.
Archaeological finds
The metro system's construction has resulted in more than 20 thousand archeological finds, from various time periods in the history of the indigenous people. The excavations needed to make way for the rails gave opportunities to find artifacts from different periods of the region's inhabitants, in areas that are now densely urbanized. Objects and small structures were found, with origins spanning from prehistoric times to the 20th century. Some examples of artifacts preserved by the Mexican National Institute of Anthropology and History (Instituto Nacional de Antropología e Historia de México (INAH)) are: parts of pyramids (like an altar to the Mexica god Ehecatl), a sculpture of the goddess Coatlicue, and remains of a mammoth. The altar to Ehécatl is now in Pino Suárez station, between lines 1 and 2, and is called by the INAH the smallest archeological site in Mexico. The metro has led to a large quantity of archeological finds, and has also let archaeologists understand more about the pattern of ancient civilisations in the Mexican capital by analysing its underground from various time periods.
Architecture
Distinguished architects were hired to design and construct the stations on the first metro line, such as Enrique del Moral, Félix Candela, Salvador Ortega and Luis Barragán. Examples of Candela's work can be seen in San Lázaro, Candelaria, and Merced stations on Line 1.
Cultural references
The Metro has figured in Mexico's cultural history, as the inspiration for a musical composition for strings, "Metro Chabacano" and Rodrigo "Rockdrigo" González's 1982 song, "Metro Balderas". It was also a filming location for the 1990 Hollywood movie Total Recall. Public intellectual Carlos Monsiváis has commented on the cultural importance of the Metro, "a space for collective expression, where diverse social sectors are compelled to mingle every day".
Lines, stations, names, colors, and logos
Each line offers one service only, and to each line, a number (letter if feeding line) and color are assigned. Every assigned color is present on square-shaped station logos, system maps and street signs, and neither colors nor numbers/letters have been changed. Line B is the only exception to the color assignment, as green (upper half) and gray (lower half) are used, producing thus bicolor logos and signs. Gray only may be used to avoid confusion with line 8, which uses a similar green.
The names of metro stations are often historical in nature, highlighting people, places, and events in Mexican history. There are stations commemorating aspects of the Mexican Revolution and the revolutionary era. When it opened in 1969 with line 1 (the "Pink Line"), two stations alluded to the Revolution. Most directly referencing the Revolution was Pino Suárez, named after Francisco I. Madero's vice president, who was murdered with him in February 1913. The other was Balderas, whose icon is a cannon, alluding to the Ciudadela armory where the coup against Madero was launched. In 1970, Revolución opened, with the station at the Monument to the Revolution. As the Metro expanded, further stations with names from the revolutionary era opened. In 1980, two popular heroes of the Revolution were honored, with Zapata explicitly commemorating the peasant revolutionary from Morelos. A sideways commemoration was División del Norte, named after the Army that Pancho Villa commanded until its demise in the Battle of Celaya in 1915.
The year 1987 saw the opening of the Lázaro Cárdenas station. In 1988, Aquiles Sedán honors the first martyr of the Revolution. In 1994, Constitución de 1917 opened, as did Garibaldi, named after the grandson of Italian fighter for independence, Giuseppe Garibaldi. The grandson had been a participant in the Mexican Revolution. In 1999, the radical anarchist Ricardo Flores Magón was honored with the station of the same name. Also opening in 1999 was Romero Rubio, named after the leader of Porfirio Díaz's Científicos, whose daughter, Carmen Romero Rubio, became Díaz's second wife. In 2012, a new Metro line opened with an Hospital 20 de Noviembre stop, a hospital named after the date that Francisco I. Madero in his 1910 Plan de San Luis Potosí called for rebellion against Díaz. There are no Metro stops named for Madero, Carranza, Obregón, or Calles, and only an oblique reference to Villa in Metro División del Norte.
Each station is identified by a minimalist logo, first designed by Lance Wyman, who had also designed the logo for the 1968 Mexico Olympics. Logos are generally related to the name of the station or the area around it. At the time of Line 1's opening, Mexico's illiteracy rate was high. As of 1960, 38% of Mexicans over the age of five were illiterate and only 5.6% of Mexicans had completed elementary school. Since one-third of the Mexican population could not read or write and most of the rest had not completed high school, it was thought that patrons would find it easier to guide themselves with a system based on colors and visual signs.
The logos are not assigned at random; rather, they are designated by considering the surrounding areas, such as:
The reference places that are located around the stations (e.g., the logo for Salto del Agua fountain depicts a fountain).
The topology of an area (e.g., Coyoacán—in Nahuatl "place of coyotes"—depicts a coyote).
The history of the place (e.g., Juárez, named after President Benito Juárez, depicts his silhouette).
The logos' background colors reflect those of the line the station serves. Stations serving two or more lines show the respective colors of each line in diagonal stripes, as in Salto del Agua. This system was adopted for the Guadalajara and Monterrey metros, and for the Mexico City Metrobús. Although logos are no longer necessary due to literacy being now widespread, their usage has remained.
Under construction:
Transfers to other systems
The Mexico City Metro offers in and out-street transfers to four major rapid transit systems: the Mexico City Metrobús and State of Mexico Mexibús bus rapid transit systems, the Mexico City light rail system and the Ferrocarril Suburbano (FSZMVM) commuter rail. None of these are part of the Sistema de Transporte Colectivo network and an extra fare must be paid for access.
Metrobús line 1 was inaugurated in 2005. According to the 1985 STC Metro Master Plan, Metrobús Line 1 roughly follows the route planned for STC Metro Line 15 by 2010, which was never built. Every transfer is out-of-station, but the same smart card may be used for payment. All five lines (Line 5 to be built during 2013) offer a connection to at least one STC Metro station. STC Metro stations that connect to Metrobús lines include Indios Verdes, La Raza, Chilpancingo, Balderas, Etiopía / Plaza de la Transparencia, Insurgentes Sur and others.
The sole light rail line running from Tasqueña to Xochimilco is operated by the Servicio de Transportes Eléctricos and is better known as Tren Ligero. Line 2 terminal Tasqueña offers an in-station transfer, but an extra ticket must be purchased.
In 2008, the Ferrocarril Suburbano commuter rail, commonly known as Suburbano, was inaugurated with a sole line running from Cuatitlán to Buenavista as of 2013. STC Metro offers two in-station transfers: Line B terminal Buenavista to the Suburbano terminal of the same name, and Line 6 station Ferrería / Arena Ciudad de México into Suburbano station Fortuna. An extra fare must be paid, and a Ferrocarril Suburbano smart card is required for access.
Another commuter rail, Tren Interurbano de Pasajeros Toluca-Valle de México is estimated to be completed in 2023. This line will connect Observatorio station in Mexico City with Toluca.
Fares and pay systems
Previously, a single ticket of MXN $5.00, allowed a rider one trip anywhere within the system with unlimited transfers. A discounted rate of MXN $3.00 is available upon application for women head of households, the unemployed, and students with scarce resources.
Mexico City Metro offers free service to the elderly, the physically impaired, and children under the age of 5 (accompanied by an adult). Tickets could be purchased at booths. They were made of paper and had a magnetic strip on them, and were recycled upon being inserted into a turnstile. As of February 2024, tickets have been discontinued and riders must obtain a rechargeable card.
Until 2009, a STC Metro ticket cost MXN $2.00 (€ 0.10, or US$ 0.15 in 2009); one purchased ticket allowed unlimited distance travel and transfer at any given time for one day, making the Mexico City Metro one of the cheapest rail systems in the world. Only line A's transfer in Pantitlán required a second payment before 13 December 2013. In January 2010, the price rose to MXN $3.00 (€ 0.15, or US$ 0.24), a fare that remained until 13 December 2013; a 2009 survey showed that 93% of citizens approved of the increase, while some said they would be willing to pay even more if needed.
STC Metro rechargeable cards were first available for an initial cost of MXN $10.00. The card would be recharged at the ticket counter in any station (or at machines in some Metro stations) to a maximum of MXN $120.00 (around € 6.44, or US$ 7.05 in 2015) for 24 trips.
In an attempt to modernize public transport, in October 2012, the Mexico City government implemented the use of a prepaid fare card, or stored-value card, called Tarjeta DF (Tarjeta del Distrito Federal, literally Federal District Card) as a payment method for STC Metro, Metrobús and the city's trolleybus and light rail systems, though they are all managed by different organizations. Servicio de Transportes Eléctricos manages both the Xochimilco Light Rail line and the city's trolleybus system. Previous fare cards that were valid only on STC Metro or Metrobús remained valid for the system for which they were acquired.
Rolling stock
As of April 2012, 14 types of standard gauge rolling stock totalling a number of 355 trains running in 6-or 9-car formation are currently in use on the Mexico City Metro. Most of the stock is rapid transit type, with the exception of the Line A stock, which is light metro. Four manufacturers have provided rolling stock for the Mexico City Metro, namely the French Alstom (MP-68, NM-73, NM-79), Canadian Bombardier (FM-95A and NM-02), Spanish CAF (NM-02, FE-07, FE-10 and NM-16 and Mexican Concarril (NM-83 and FM-86) (now Bombardier Transportation Mexico, in some train types with the help of Alstom and/or Bombardier).
The maximum design speed limit is (average speed ) for rubber-tired rolling stock and (average speed ) for steel-wheeled rolling stock. Forced-air ventilation is employed and the top portion of windows can be opened so that passenger comfort is enhanced by the combination of these two types of ventilation. Like the rolling stock used in the Paris Métro and the Montreal Metro, the numbering of the Mexico City Metro's rolling stock are specified by year of design (not year of first use).
In chronological order, the types of rubber-tired rolling stock are: MP-68, NM-73A, NM-73B, NM-73C, NM-79, MP-82, NC-82, NM-83A, NM-83B, NE-92, NM-02, NM-16 and NM-22; and the types of steel-wheeled rolling stock are: FM-86, FM-95A, FE-07, and FE-10.
From May 2024, Line 1 will receive 29 new rubber-tired trains manufactured by CRRC Zhuzhou Locomotive in China, replacing earlier rolling stock. This is in line with ongoing upgrading works for Line 1, including the installation of CBTC.
Gallery
Major incidents
On 20 October 1975, two trains crashed in Viaducto station while both were going towards Tasqueña station. The first was stopped picking up passengers when it was hit by another train that did not stop in time. According to official reports, from 31 to 39 people died, and between 71 and 119 were injured. After the crash, automatic signals were incorporated to all lines.
On 18 September 2009, a man was vandalizing the walls of Balderas station with a marker before being confronted by a police officer. He took out a gun and killed the officer and a construction worker who tried to disarm him, and injured 5 others.
On 4 May 2015, two trains heading towards Politécnico station on Line 5 crashed in Oceanía station. The first was leaving to Aragón station and was requested to stop and wait, while the second did not deactivate the autopilot and crashed into it at the end of the platform. 12 people were injured.
On 10 March 2020, two trains heading towards Observatorio station on Line 1 crashed in Tacubaya station. The first train was parked at the platform when it was hit by another train that was coming in reverse. 1 person died and 41 were injured, all inside the second train, as people in the parked train had been evacuated moments before the crash.
On 9 January 2021, the Central Control Center serving lines 1 to 6 caught fire. During the fire, a female police officer was killed due to a fall in the building. All the stations on those lines temporarily remained closed and provisional transport service was provided by city buses and police vehicles. According to the Metro authorities, the service in lines 4, 5, and 6 would be normalized in days, while that in lines 1, 2, and 3 in several months.
On 3 May 2021, a train was traveling on Line 12 between the Olivos and Tezonco stations when a girder supporting the overpass on which the train was traveling collapsed, killing 26 and injuring more than 70. Service on Line 12 was later suspended, while STC warned residents to avoid the site of the collapse.
On 7 January 2023 at 09:16 local time, two trains collided between Potrero and La Raza stations on Line 3, killing one and injuring 57. In addition to other minor events, city officials said that this accident was a result of sabotage to the Fourth Transformation platform to affect the image of Claudia Sheinbaum, then mayor of the city and a potential candidate on the 2024 Mexican general election.
| Technology | Americas | null |
892484 | https://en.wikipedia.org/wiki/Anaerobic%20exercise | Anaerobic exercise | Anaerobic exercise is a type of exercise that breaks down glucose in the body without using oxygen; anaerobic means "without oxygen". This type of exercise leads to a buildup of lactic acid.
In practical terms, this means that anaerobic exercise is more intense, but shorter in duration than aerobic exercise.
The biochemistry of anaerobic exercise involves a process called glycolysis, in which glucose is converted to adenosine triphosphate (ATP), the primary source of energy for cellular reactions.
Anaerobic exercise may be used to help build endurance, muscle strength, and power.
Metabolism
Anaerobic metabolism is a natural part of metabolic energy expenditure. Fast twitch muscles (as compared to slow twitch muscles) operate using anaerobic metabolic systems, such that any use of fast twitch muscle fibers leads to increased anaerobic energy expenditure. Intense exercise lasting upwards of four minutes (e.g. a mile race) may still have considerable anaerobic energy expenditure. An example is high-intensity interval training, an exercise strategy that is performed under anaerobic conditions at intensities that reach an excess of 90% of the maximum heart rate. Anaerobic energy expenditure is difficult to accurately quantify. Some methods estimate the anaerobic component of an exercise by determining the maximum accumulated oxygen deficit or measuring the lactic acid formation in muscle mass.
In contrast, aerobic exercise includes lower intensity activities performed for longer periods of time. Activities such as walking, jogging, rowing, and cycling require oxygen to generate the energy needed for prolonged exercise (i.e., aerobic energy expenditure). For sports that require repeated short bursts of exercise, the aerobic system acts to replenish and store energy during recovery periods to fuel the next energy burst. Therefore, training strategies for many sports demand that both aerobic and anaerobic systems be developed. The benefits of adding anaerobic exercise include improving cardiovascular endurance as well as build and maintaining muscle strength and losing weight.
The anaerobic energy systems are:
The alactic anaerobic system, which consists of high energy phosphates, adenosine triphosphate, and creatine phosphate; and
The lactic anaerobic system, which features anaerobic glycolysis.
High energy phosphates are stored in limited quantities within muscle cells. Anaerobic glycolysis exclusively uses glucose (and glycogen) as a fuel in the absence of oxygen, or more specifically, when ATP is needed at rates that exceed those provided by aerobic metabolism. The consequence of such rapid glucose breakdown is the formation of lactic acid (or more appropriately, its conjugate base lactate at biological pH levels). Physical activities that last up to about thirty seconds rely primarily on the former ATP-CP phosphagen system. Beyond this time, both aerobic and anaerobic glycolysis-based metabolic systems are used.
The by-product of anaerobic glycolysis—lactate—has traditionally been thought to be detrimental to muscle function. However, this appears likely only when lactate levels are very high. Elevated lactate levels are only one of many changes that occur within and around muscle cells during intense exercise that can lead to fatigue. Fatigue, which is muscle failure, is a complex subject that depends on more than just changes to lactate concentration. Energy availability, oxygen delivery, perception to pain, and other psychological factors all contribute to muscular fatigue. Elevated muscle and blood lactate concentrations are a natural consequence of any physical exertion. The effectiveness of anaerobic activity can be improved through training.
Anaerobic exercise also increases an individual's basal metabolic rate (BMR).
Examples
Anaerobic exercises are high-intensity workouts completed over shorter durations, while aerobic exercises include variable-intensity workouts completed over longer durations. Some examples of anaerobic exercises include sprints, high-intensity interval training (HIIT), and strength training.
| Biology and health sciences | Physical fitness | Health |
892541 | https://en.wikipedia.org/wiki/Dicynodontia | Dicynodontia | Dicynodontia is an extinct clade of anomodonts, an extinct type of non-mammalian therapsid. Dicynodonts were herbivores that typically bore a pair of tusks, hence their name, which means 'two dog tooth'. Members of the group possessed a horny, typically toothless beak, unique amongst all synapsids. Dicynodonts first appeared in Southern Pangaea during the mid-Permian, ca. 270–260 million years ago, and became globally distributed and the dominant herbivorous animals in the Late Permian, ca. 260–252 Mya. They were devastated by the end-Permian Extinction that wiped out most other therapsids ca. 252 Mya. They rebounded during the Triassic but died out towards the end of that period. They were the most successful and diverse of the non-mammalian therapsids, with over 80-90 genera known, varying from rat-sized burrowers to elephant-sized browsers.
Characteristics
The dicynodont skull is highly specialised, light but strong, with the synapsid temporal openings at the rear of the skull greatly enlarged to accommodate larger jaw muscles. The front of the skull and the lower jaw are generally narrow and, in all but a number of primitive forms, toothless. Instead, the front of the mouth is equipped with a horny beak, as in turtles and ceratopsian dinosaurs. Food was processed by the retraction of the lower jaw when the mouth closed, producing a powerful shearing action, which would have enabled dicynodonts to cope with tough plant material. Dicynodonts typically had a pair of enlarged maxillary caniniform teeth, analogous to the tusks present in some living mammals. In the earliest genera, they were merely enlarged teeth, but in later forms they independently evolved into ever-growing teeth like mammal tusks multiple times. In some dicynodonts, the presence of tusks has been suggested to be sexually dimorphic. Some dicynodonts such as Stahleckeria lacked true tusks and instead bore tusk-like extensions on the side of the beak.
The body is short, strong and barrel-shaped, with strong limbs. In large genera (such as Dinodontosaurus) the hindlimbs were held erect, but the forelimbs bent at the elbow. Both the pectoral girdle and the ilium are large and strong. The tail is short.
Pentasauropus dicynodont tracks suggest that dicynodonts had fleshy pads on their feet. Mummified skin from specimens of Lystrosaurus in South Africa have numerous raised bumps.
Endothermy and soft tissue anatomy
Dicynodonts have long been suspected of being warm-blooded animals. Their bones are highly vascularised and possess Haversian canals, and their bodily proportions are conducive to heat preservation. In young specimens, the bones are so highly vascularised that they exhibit higher channel densities than most other therapsids. Yet, studies on Late Triassic dicynodont coprolites paradoxically showcase digestive patterns more typical of animals with slow metabolisms.
More recently, the discovery of hair remnants in Permian coprolites possibly vindicates the status of dicynodonts as endothermic animals. As these coprolites come from carnivorous species and digested dicynodont bones are abundant, it has been suggested that at least some of these hair remnants come from dicynodont prey. A new study using chemical analysis seemed to suggest that cynodonts and dicynodonts both developed warm blood independently before the Permian extinction.
History
A 2024 paper posited that rock art of a superficially walrus-like imaginary creature with downcurved tusks created by the San people of South Africa prior to 1835 may have been partly inspired by fossil dicynodont skulls which erode out of rocks in the area.
Dicynodonts have been known to science since the mid-1800s. The South African geologist Andrew Geddes Bain gave the first description of dicynodonts in 1845. At the time, Bain was a supervisor for the construction of military roads under the Corps of Royal Engineers and had found many reptilian fossils during his surveys of South Africa. Bain described these fossils in an 1845 letter published in Transactions of the Geological Society of London, calling them "bidentals" for their two prominent tusks. In that same year, the English paleontologist Richard Owen named two species of dicynodonts from South Africa: Dicynodon lacerticeps and Dicynodon bainii. Since Bain was preoccupied with the Corps of Royal Engineers, he wanted Owen to describe his fossils more extensively. Owen did not publish a description until 1876 in his Descriptive and Illustrated Catalogue of the Fossil Reptilia of South Africa in the Collection of the British Museum. By this time, many more dicynodonts had been described. In 1859, another important species called Ptychognathus declivis was named from South Africa. In the same year, Owen named the group Dicynodontia. In his Descriptive and Illustrated Catalogue, Owen honored Bain by erecting Bidentalia as a replacement name for his Dicynodontia. The name Bidentalia quickly fell out of use in the following years, replaced by popularity of Owen's Dicynodontia.
Evolutionary history
Dicynodonts first appeared during the Middle Permian in the Southern Hemisphere, with South Africa being the centre of their known diversity, and underwent a rapid evolutionary radiation, becoming globally distributed and amongst the most successful and abundant land vertebrates during the Late Permian. During this time, they included a large variety of ecotypes, including large, medium-sized, and small herbivores and short-limbed mole-like burrowers.
Only four lineages are known to have survived the Great Dying; the first three represented with a single genus each: Myosaurus, Kombuisia, and Lystrosaurus, the latter being the most common and widespread herbivores of the Induan (earliest Triassic). None of these survived long into the Triassic. The fourth group was the Kannemeyeriiformes, the only dicynodonts who diversified during the Triassic. These stocky, pig- to ox-sized animals were the most abundant herbivores worldwide from the Olenekian to the Ladinian age. By the Carnian they had been supplanted by traversodont cynodonts and rhynchosaur reptiles. During the Norian (middle of the Late Triassic), perhaps due to increasing aridity, they drastically declined, and the role of large herbivores was taken over by sauropodomorph dinosaurs.
Fossils of an Asian elephant-sized dicynodont Lisowicia bojani discovered in Poland indicate that dicynodonts survived at least until the late Norian or earliest Rhaetian (latest Triassic); this animal was also the largest known dicynodont species.
Six fragments of fossil bone discovered in Queensland, Australia, were interpreted as remains of a skull in 2003. This suggested to indicate that dicynodonts survived into the Cretaceous in southern Gondwana. The dicynodont affinity of these specimens was questioned (including a proposal that they belonged to a baurusuchian crocodyliform by Agnolin et al. in 2010), and in 2019 Knutsen and Oerlemans considered this fossil to be of Plio-Pleistocene age, and reinterpreted it as a fossil of a large mammal, probably a diprotodontid.
With the decline and extinction of the kannemeyerids, there were to be no more dominant large synapsid herbivores until the middle Paleocene epoch (60 Ma) when mammals, distant descendants of cynodonts, began to diversify after the extinction of the non-avian dinosaurs.
Systematics
Taxonomy
Dicynodontia was originally named by the English paleontologist Richard Owen. It was erected as a family of the order Anomodontia and included the genera Dicynodon and Ptychognathus. Other groups of Anomodontia included Gnathodontia, which included Rhynchosaurus (now known to be an archosauromorph) and Cryptodontia, which included Oudenodon. Cryptodonts were distinguished from dicynodonts from their absence of tusks. Although it lacks tusks, Oudenodon is now classified as a dicynodont, and the name Cryptodontia is no longer used. Thomas Henry Huxley revised Owen's Dicynodontia as an order that included Dicynodon and Oudenodon. Dicynodontia was later ranked as a suborder or infraorder with the larger group Anomodontia, which is classified as an order. The ranking of Dicynodontia has varied in recent studies, with Ivakhnenko (2008) considering it a suborder, Ivanchnenko (2008) considering it an infraorder, and Kurkin (2010) considering it an order.
Many higher taxa, including infraorders and families, have been erected as a means of classifying the large number of dicynodont species. Cluver and King (1983) recognised several main groups within Dicynodontia, including Eodicynodontia (containing only Eodicynodon), Endothiodontia (containing only Endothiodontidae), Pristerodontia (Pristerodontidae, Cryptodontidae, Aulacephalodontidae, Dicynodontidae, Lystrosauridae, and Kannemeyeriidae), Kingoriamorpha (containing only Kingoriidae), Diictodontia (Diictodontidae, Robertiidae, Cistecephalidae, Emydopidae and Myosauridae), and Venyukoviamorpha. Most of these taxa are no longer considered valid. Kammerer and Angielczyk (2009) suggested that the problematic taxonomy and nomenclature of Dicynodontia and other groups results from the large number of conflicting studies and the tendency for invalid names to be mistakenly established.
Phylogeny
Below is a cladogram modified from Angielczyk et al. (2021):
Current classification
Dicynodontia
Brachyprosopus
Colobodectes
Eodicynodon
Lanthanostegus
Nyaphulia
Endothiodontia
Abajudon
Endothiodon
Niassodon
Eumantellidae
Pristerodon
Pylaecephalidae
Diictodon
Eosimops
Prosictodon
Robertia
Therochelonia
Emydopoidea
Digalodon
Rastodon
Cistecephalidae
Cistecephalus
Cistecephaloides
Kawingasaurus
Kembawacela
Sauroscaptor
Emydopidae
Compsodon
Emydops
Kingoriidae
Dicynodontoides
Kombuisia
Thliptosaurus
Myosauridae
Myosauroides
Myosaurus
Bidentalia
Kunpania
Cryptodontia
Daqingshanodon
Keyseria
Rhachiocephalidae
Kitchinganomodon
Rhachiocephalus
Oudenodontidae
Australobarbarus
Oudenodon
Tropidostoma
Geikiidae
Bulbasaurus
?Idelesaurus
?Odontocyclops
Geikiinae
Aulacephalodon
Geikia
Pelanomodon
Dicynodontoidea
Counillonia
Daptocephalus
Delectosaurus
Dicynodon
Dinanomodon
Elph
Gordonia
Interpresosaurus
Katumbia
Peramodon
Taoheodon
Turfanodon
Vivaxosaurus
Lystrosauridae
?Basilodon
?Jimusaria
?Sintocephalus
?Syops
Euptychognathus
Kwazulusaurus
Lystrosaurus
Kannemeyeriiformes
Angonisaurus
Dinodontosauridae
Dinodontosaurus
Shansiodontidae
Rhinodicynodon
Shansiodon
Tetragonias
Vinceria
Kannemeyeriidae
Acratophorus
Dolichuranus
Kannemeyeria
Parakannemeyeria
Rabidosaurus
Rechnisaurus
Rhadiodromus
Shaanbeikannemeyeria
Sinokannemeyeria
Uralokannemeyeria
Wadiasaurus
Xiyukannemeyeria
Stahleckeriidae
?Sungeodon
Woznikella
Placeriinae
Argodicynodon
Lisowicia
Moghreberia
Placerias
Pentasaurus
Zambiasaurus
Stahleckeriinae
Eubrachiosaurus
Ischigualastia
Jachaleria
Sangusaurus
Stahleckeria
Ufudocyclops
South African geomyth
A horned serpent cave art is known from the La Belle France cave in South Africa, often conflated with the Dingonek. It may be based on dicynodont fossils.
| Biology and health sciences | Proto-mammals | Animals |
892554 | https://en.wikipedia.org/wiki/Placerias | Placerias | Placerias (meaning 'broad body') is an extinct genus of dicynodonts that lived during the Carnian to the Norian age of the Triassic Period (230–215 million years ago). Placerias belongs to a group of dicynodonts called Kannemeyeriiformes, which was the last known group of dicynodonts before the taxon became extinct at the end of the Triassic.
Description
Placerias was one of the largest herbivores in the Late Triassic, weighing up to . The largest skull found had a length of .
Placerias had a powerful neck, strong legs, and barrel-shaped body with possible ecological and evolutionary parallels with the modern hippopotamus, spending much of its time during the wet season wallowing in the water and chewing at bankside vegetation. Placerias was closely related to Ischigualastia and similar in appearance. Placerias used its beak to slice through thick branches and roots with two short tusks that could be used for defence and for intra-specific display. The genus exhibits two morphs, one with short tusks and one with long tusks, which is inferred to be sexual dimorphism, with the longer-tusked individuals presumably being males.
Discovery
Fossils of forty Placerias were found near St. Johns, southeast of the Petrified Forest in the Chinle Formation of Arizona. This site has become known as the 'Placerias Quarry' and was discovered in 1930, by Charles Camp and Samuel Welles, of the University of California, Berkeley. Sedimentological features of the site indicate a low-energy depositional environment, possibly flood-plain or overbank. Bones are associated mostly with mudstones and a layer that contains numerous carbonate nodules. It is also known from the Pekin Formation of North Carolina.
Placerias was originally considered the last of the dicynodonts, although other Late Triassic dicynodonts, such as Lisowicia and Pentasaurus have since been discovered.
| Biology and health sciences | Proto-mammals | Animals |
892600 | https://en.wikipedia.org/wiki/Eurypterus | Eurypterus | Eurypterus ( ) is an extinct genus of eurypterid, a group of organisms commonly called "sea scorpions". The genus lived during the Silurian period, from around 432 to 418 million years ago. Eurypterus is by far the most well-studied and well-known eurypterid. Eurypterus fossil specimens probably represent more than 95% of all known eurypterid specimens.
There are fifteen species belonging to the genus Eurypterus, the most common of which is E. remipes, the first eurypterid fossil discovered and the state fossil of New York.
Members of Eurypterus averaged at about in length, but the largest individual discovered was estimated to be long. They all possessed spine-bearing appendages and a large paddle they used for swimming. They were generalist species, equally likely to engage in predation or scavenging.
Discovery
The first fossil of Eurypterus was found in 1818 by S. L. Mitchill, a fossil collector. It was recovered from the Bertie Formation of New York (near Westmoreland, Oneida County). Mitchill interpreted the appendages on the carapace as barbels arising from the mouth. He consequently identified the fossil as a catfish of the genus Silurus.
It was only after seven years, in 1825, that the American zoologist James Ellsworth De Kay identified the fossil correctly as an arthropod. He named it Eurypterus remipes and established the genus Eurypterus in the process. The name means 'wide wing' or 'broad paddle', referring to the swimming legs, from Greek ( 'wide') and ( 'wing').
However, De Kay thought Eurypterus was a branchiopod (a group of crustaceans which include fairy shrimps and water fleas). Soon after, Eurypterus lacustris was also discovered in New York in 1835 by the paleontologist Richard Harlan. Another species was discovered in Estonia in 1858 by Jan Nieszkowski. He considered it to be of the same species as the first discovery (E. remipes); it is now known as E. tetragonophthalmus. These specimens from Estonia are often of extraordinary quality, retaining the actual cuticle of their exoskeletons. In 1898, the Swedish paleontologist Gerhard Holm separated these fossils from the bedrock with acids. Holm was then able to examine the almost perfectly preserved fragments under a microscope. His remarkable study led to the modern breakthrough on eurypterid morphology.
More fossils were recovered in great abundance in New York in the 19th century, and elsewhere in eastern Eurasia and North America. Today, Eurypterus remains one of the most commonly found and best known eurypterid genera, comprising more than 95% of all known eurypterid fossils.
E. remipes was designated the New York State Fossil by the then Governor Mario Cuomo in 1984.
Description
The largest arthropods to have ever existed were eurypterids. The largest known species (Jaekelopterus rhenaniae) reached up to in length, about the size of a crocodile. Species of Eurypterus, however, were much smaller.
E. remipes are usually between in length. E. lacustris average at larger sizes at in length. However, a single telson (the posteriormost division of the body) of a specimen of this species reaches this length, being long and indicating a specimen of of length, and that is the largest specimen ever described in literature. In the introduction page of E. remipes in website of University of Texas at Austin says that the largest specimen ever found was long, currently on display at the Paleontological Research Institution of New York. However, the text section describes the group eurypterid itself rather than Eurypterus, so it is not possible to determine in context whether the long specimen is actually from E. remipes or another eurypterid.
Eurypterus fossils often occur in similar sizes in a given area. This may be a result of the fossils being "sorted" into windrows as they were being deposited in shallow waters by storms and wave action.
The Eurypterus body is broadly divided into two parts: the prosoma and the opisthosoma (in turn divided into the mesosoma and the metasoma).
The prosoma is the forward part of the body, it is actually composed of six segments fused together to form the head and the thorax. It contains the semicircular to subrectangular platelike carapace. On the dorsal side of the latter are two large crescent-shaped compound eyes. They also possessed two smaller light-sensitive simple eyes (the median ocelli) near the center of the carapace on a small elevation (known as the ocellar mound). Underneath the carapace is the mouth and six appendages, usually referred to in Roman numerals I-VI. Each appendage in turn is composed of nine segments (known as podomeres) labeled in Arabic numerals 1–9. The first segments which connect the appendages to the body are known as the coxa (plural coxae).
The first pair (Appendage I) are the chelicerae, small pincer-like arms used for tearing food apart (mastication) during feeding. After the chelicerae are three pairs of short legs (Appendages II, III, and IV). They are spiniferous, with predominantly two spines on each podomere and with the tipmost segment having a single spine. The last two segments are often indistinguishable and give the appearance of a single segment having three spines. They are used both for walking and for food capture. The next pair (Appendage V) is the most leg-like of all appendages, longer than the first three pairs and are mostly spineless except at the tipmost segments. The last pair (Appendage VI) are two broad paddle-like legs used for swimming. The coxae of Appendage VI are broad and flat, resembling an 'ear'.
The ophisthosoma (the abdomen) is composed of 12 segments, each consisting of a fused upper plate (tergite) and bottom plate (sternite). It is further subdivided in two ways.
Based on the width and structure of each segment, they can be divided into the broad preabdomen (segments 1 to 7) and the narrow postabdomen (segments 8 to 12). The preabdomen is the broader segments of the anterior portion of the ophisthosoma while the postabdomen are the last five segments of the Eurypterus body. Each of the segments of the postabdomen contain lateral flattened protrusions known as the epimera with the exception of the last needle-like (styliform) part of the body known as the telson. The segment immediately preceding the telson (which also has the largest epimera of the postabdomen) is known as the pretelson.
An alternative way to divide the ophisthosoma is by function. It can also be divided into the mesosoma (segments 1 to 6), and the metasoma (segments 7 to 12). The mesosoma contains the gills and reproductive organs of Eurypterus. Its ventral segments are overlaid by appendage-derived plates known as Blattfüsse (singular Blattfuss, German for "sheet foot"). Protected within which are the branchial chambers which contain the respiratory organs of Eurypterus. The metasoma, meanwhile, do not possess Blattfüsse.
Some authors incorrectly use mesosoma and preabdomen interchangeably, as with metasoma and postabdomen.
The main respiratory organs of Eurypterus were what seems to be book gills, located in branchial chambers within the segments of the mesosoma. They may have been used for underwater respiration. They are composed of several layers of thin tissue stacked in such a way as to resemble the pages of a book, hence the name. In addition, they also possessed five pairs of oval-shaped areas covered with microscopic projections on the ceiling of the second branchial chambers within the mesosoma, immediately below the gill tracts. These areas are known as Kiemenplatten (or gill-tracts, though the former term is preferred). They are unique to eurypterids.
Eurypterus are sexually dimorphic. On the bottom side of the first two segments of the mesosoma are central appendages used for reproduction. In females, they are long and narrow. In the males they are very short. A minority of authors, however, assume the reverse: longer genital appendage for males, shorter for females.
The exoskeleton of Eurypterus is often covered with small outgrowths known as ornamentation. They include pustules (small protrusions), scales, and striations. They vary by species and are used for identification. For more detailed diagnostic descriptions of each species under Eurypterus, see sections below.
Classification
The genus Eurypterus belongs to the family Eurypteridae. They are classified under the superfamily Eurypteroidea, suborder Eurypterina, order Eurypterida, and the subphylum Chelicerata. Until recently, eurypterids were thought to belong to the class Merostomata along with order Xiphosura. It is now believed that eurypterids are a sister group to Arachnida, closer to scorpions and spiders than to horseshoe crabs.
Eurypterus was the first recognized taxon of eurypterids and is the most common. As a consequence, nearly every remotely similar eurypterid in the 19th century was classified under the genus (except for the distinctive members of the family Pterygotidae and Stylonuridae). The genus was eventually split into several genera as the science of taxonomy developed.
In 1958, several species distinguishable by closer placed eyes and spines on their swimming legs were split off into the separate genus Erieopterus by Erik Kjellesvig-Waering. Another split was proposed by Leif Størmer in 1973 when he reclassified some Eurypterus to Baltoeurypterus based on the size of some of the last segments of their swimming legs. O. Erik Tetlie in 2006 deemed these differences too insignificant to justify a separate genus. He merged Baltoeurypterus back into Eurypterus. It is now believed that the minor variations described by Størmer are simply the differences found in adults and juveniles within a species.
The genus Eurypterus derives from E. minor, the oldest known species from the Llandovery of Scotland. E. minor is believed to have diverged from Dolichopterus macrocheirus sometime in the Llandovery. The following is the phylogenetic tree of Eurypterus based on phylogenetic studies by O. Erik Tetlie in 2006. Some species are not represented.
Species
Species belonging to the genus, their diagnostic descriptions, synonyms (if present), and distribution are as follows:
Eurypterus De Kay, 1825
?Eurypterus cephalaspis Salter, 1856 – Silurian, England
Uncertain placement. Only 3 of the specimens described in 1856 are probably Eurypterus, the rest probably belonged to Hughmilleriidae. Its name means 'shield head', from Greek ( 'head'), and ( 'shield or bowl'). Specimens recovered from Herefordshire, England.
Eurypterus dekayi Hall, 1859 – Silurian, United States & Canada
No raised scales on the posterior margin of the carapace or of the three front-most tergites. The rest of the tergites each have four raised scales. Four to six spines on each podomere of Appendages III and IV. Pretelson has large, rounded epimera without ornamentation on the margins. The species is very similar to E. laculatus. The species is named after James Ellsworth De Kay. Specimens recovered from New York and Ontario.
Eurypterus flintstonensis Swartz, 1923 – Silurian, USA
Probably a synonym of E. remipes or E. lacustris. Probably named after Flintstone, Georgia (?). Specimen recovered from eastern United States.
Eurypterus hankeni Tetlie, 2006 – Silurian, Norway
Small Eurypterus species, averaging at long. The largest specimen found is about in length. They can be distinguished by pustules and six scales at the rear margin of their carapaces. Appendages I to IV has two spines on each podomere. The postabdomen have small epimera. The pretelson has long pointed epimera. Telson has striations near its attachment to the pretelson. The species is named after Norwegian paleontologist Nils-Martin Hanken, of the University of Tromsø. Found in the Steinsfjorden Formation of Ringerike, Norway.
Eurypterus henningsmoeni Tetlie, 2002 – Silurian, Norway
Eurypterus with broad paddles and metastoma. Postabdomen has small epimera. Pretelson has large rounded epimera with imbricate scales (overlapping, similar to fish scales). It is very similar and closely related to E. tetragonophthalmus. The species was named after the Norwegian paleontologist Gunnar Henningsmoen. Found in Bærum, Norway.
Eurypterus laculatus Kjellesvig-Waering, 1958 – Silurian, USA & Canada
The visual area of the compound eyes of this species are surrounded by depressions. The ocelli and the ocellar mound are small. No pustules or raised scales on the carapace or the first tergite. It is probably closely related to E. dekayi. Its specific epithet means 'four-cornered', from Latin ('four-cornered, checkered'). Found in New York and Ontario.
Eurypterus lacustris Harlan, 1834 – Silurian, USA & Canada
= Eurypterus pachycheirus Hall, 1859 – Silurian, USA & Canada
= Eurypterus robustus Hall, 1859 – Silurian, USA & Canada
One of the two most common Eurypterus fossils found. It is very similar to E. remipes and often found in the same localities, but the eyes are placed at a more posterior position on the carapace of E. lacustris. It is also slightly larger with a slightly narrower metastoma. Its status as a distinct species was once disputed before diagnostic analysis by Tollerton in 1993. Its specific name means 'from a lake', from Latin ('lake'). Found in New York and Ontario.
Eurypterus leopoldi Tetlie, 2006 – Silurian, Canada
Frontmost tergite is reduced. Metasoma is rhombiovate in shape with tooth-like projections at the anterior part. The pretelson has serrated edges. the epimera are large, semi-angular with angular striations. The telson is styliform with large angular striations interspersed among smaller more numerous striations. The species is named after Port Leopold and the Leopold Formation where they were collected. Found in the Leopold Formation of Somerset Island, Canada.
Eurypterus megalops Clarke & Ruedemann, 1912 – Silurian, USA
Specific name means "large eye", from Greek ( 'big or large') and ( 'eye'). Discovered in New York, United States.
?Eurypterus minor Laurie, 1899 – Silurian, Scotland
Small Eurypterus with large pustules on the carapace and abdomen. Does not possess the scale ornamentation found in other species of Eurypterus. It is the earliest known species of Eurypterus. They have large palpebral lobes (part of "cheeks" of the carapace adjacent to the compound eyes), making it easy to mistake their eyes for being oval. This enlargement is more typical of the genus Dolichopterus and it may actually belong to Dolichopteridae. The specific name means 'smaller', from Latin . Found in the Reservoir Formation of Pentland Hills, Scotland.
Eurypterus ornatus Leutze, 1958 – Silurian, USA
Ornamentation of pustules on the entire surface of the carapace and at least the first tergite. Does not possess raised scales. Its specific name means 'adorned', from Latin ('adorned, ornate'). Recovered from Fayette, Ohio.
Eurypterus pittsfordensis Sarle, 1903 – Silurian, USA
The posterior margin of the carapace has three raised scales. Appendages II to IV has two spines per podomere. The metastoma is rhomboid in shape with a deep notch at the front part. The postabdomen has serrated fringes at the middle with small angular epimera at the sides. The pretelson has large, semiangular epimera with angular striations at the margins. The telson is styliform with sparse angular striations at the margins. The name of the species comes from its place of discovery – the Salina shale formations of Pittsford, New York.
Eurypterus quebecensis Kjellesvig-Waering, 1958 – Silurian, Canada
Has six raised scales on the posterior margin of the carapace but does not possess pustule ornamentation. It is named after the location it was recovered from Quebec, Canada.
Eurypterus remipes DeKay, 1825 – Silurian, USA, Canada
= Carcinosoma trigona (Ruedemann, 1916) – Silurian, USA
The most common Eurypterus species. Has four raised scales at the posterior margin of the carapace. Appendages I to IV has two spines on each podomere. Postabdomen has small epimera. Pretelson has small, semiangular epimera with imbricate scale ornamentation at the margins. The telson has serrated margins along most of its length. It is very similar to E. lacustris and can often only be distinguished by the position of the eyes. The specific name means 'oar-foot', from Latin ('oar') and ('foot'). Found in New York and Ontario, and is the state fossil of New York.
Eurypterus serratus (Jones & Woodward, 1888) – Silurian, Sweden
Similar to E. pittsfordensis and E. leopoldi but can be distinguished by the dense angular striations on their styliform telson. The specific name means 'serrated', from Latin ('sawn [into pieces]'). Originally discovered from Gotland, Sweden.
Eurypterus tetragonophthalmus Fischer, 1839 – Silurian, Ukraine & Estonia
= Eurypterus fischeri Eichwald, 1854 – Silurian, Ukraine
= Eurypterus fischeri var. rectangularis Schmidt, 1883 – Silurian, Estonia
Four raised scales on the posterior margin of the carapce. Appendages II to IV each have two spines on each podomere. Postabdomen has small epimera. The pretelson has large, rounded epimera with imbricate scale ornamentation at the margins. Telson has imbricate scale ornamentations at the margins of the base which become serrations towards the tip. The specific name means 'four-edged eye', from Greek ( 'four'), ( 'angle'), and ( 'eye'). Found in the Rootsiküla Formation of Saaremaa (Ösel), Estonia, with additional discoveries in Ukraine, Norway, and possibly Moldova and Romania.
The list does not include the large number of fossils previously classified under Eurypterus. Most of them are now reclassified to other genera, identified as other animals (like crustaceans) or pseudofossils, or remains of doubtful placement. Classification is based on Dunlop et al.(2011).
Paleobiology
Eurypterus belongs to the suborder Eurypterina, eurypterids in which the sixth appendage had developed a broad swimming paddle remarkably similar to that of the modern-day swimming crab. Modeling studies on Eurypterus swimming behavior suggest that they utilized a drag-based rowing type of locomotion where appendages moved synchronously in near-horizontal planes. The paddle blades are almost vertically oriented on the backward and down stroke, pushing the animal forward and lifting it up. The blades are then oriented horizontally on the recovery stroke to slash through the water without pushing the animal back. This type of swimming is exhibited by crabs and water beetles.
An alternative hypothesis for Eurypterus swimming behavior is that individuals were capable of underwater flying (or subaqueous flight), in which the sinuous motions and shape of the paddles themselves acting as hydrofoils are enough to generate lift. This type is similar to that found in sea turtles and sea lions. It has a relatively slower acceleration rate than the rowing type, especially since adults have proportionally smaller paddles than juveniles. But since the larger sizes of adults mean a higher drag coefficient, using this type of propulsion is more energy-efficient.
Juveniles probably swam using the rowing type, the rapid acceleration afforded by this propulsion is more suited for quickly escaping predators. A small Eurypterus could achieve two and a half body lengths per second immediately. Larger adults, meanwhile, probably swam with the subaqueous flight type. The maximum velocity of adults when cruising would have been per second, slightly faster than turtles and sea otters.
Trace fossil evidence indicates that Eurypterus employed a rowing stroke when in close proximity to the seafloor. Arcuites bertiensis is an ichnospecies that includes a pair of crescent-shaped impressions and a short medial drag, and it has been found in upper Silurian eurypterid Lagerstatten in Ontario and Pennsylvania. This trace fossil is very similar to traces made by modern aquatic swimming insects that row such as water boatmen, and is considered to have been made by juvenile to adult-sized eurypterids while swimming in very shallow nearshore marine environments. The morphology of A. bertiensis suggests that Eurypterus had the ability to move its swimming appendages in both the horizontal and vertical plane.
Eurypterus did not swim to hunt, rather they simply swam in order to move from one feeding site to another quickly. Most of the time they walked on the substrate with their legs (including their swimming leg). They were generalist species, equally likely to engage in predation or scavenging. They hunted small soft-bodied invertebrates like worms. They utilized the mass of spines on their front appendages to both kill and hold them while they used their chelicerae to rip off pieces small enough to swallow. Young individuals may also have fallen prey to cannibalism by larger adults.
Eurypterus were most probably marine animals, as their remains are mostly found in intertidal shallow environments. The concentrations of Eurypterus fossils in certain sites has been interpreted to be a result of mass mating and molting behavior. Juveniles were likely to have inhabited nearshore hypersaline environments, safer from predators, and moved to deeper waters as they grew older and larger. Adults that reach sexual maturity would then migrate en masse to shore areas in order to mate, lay eggs, and molt. Activities that would have made them more vulnerable to predators. This could also explain why the vast majority of fossils found in such sites are molts and not of actual animals. The same behavior can be seen in modern horseshoe crabs.
Respiration
Examinations of the respiratory systems of Eurypterus have led many paleontologists to conclude that it was capable of breathing air and walking on land for a short amount of time. Eurypterus had two types of respiratory systems. Its main organs for breathing were the book gills inside the segments of the mesosoma. These structures were supported by semicircular 'ribs' and were probably attached near the center of the body, similar to the gills of modern horseshoe crabs. They were protected under platelike appendages (which actually formed the apparent 'belly' of Eurypterus) known as Blattfüsse. These gills may have also played a role in osmoregulation.
The second system are the Kiemenplatten, also referred to as gill-tracts. These oval-shaped areas within the body wall of the preabdomen. Their surfaces are covered with numerous small spines arranged into hexagonal 'rosettes'. These areas were vascularized, hence the conclusion that they were secondary breathing organs.
The function of the book gills are usually interpreted to be for aquatic breathing, while the Kiemenplatten are supplementary for temporary breathing on land. However, some authors have argued that the two systems alone could not have supported an organism the size of Eurypterus. Both structures might actually have been for breathing air and the true gills (for underwater breathing) of Eurypterus have yet to be discovered. Eurypterus, however, were undoubtedly primarily aquatic.
Ontogeny
Juvenile Eurypterus differed from adults in several ways. Their carapaces were narrower and longer (parabolic) in contrast to the trapezoidal carapaces of adults. The eyes are aligned almost laterally but move to a more anterior location during growth. The preabdomen also lengthened, increasing the overall length of the ophisthosoma. The swimming legs also became narrower and the telsons shorter and broader (though in E. tetragonophthalmus and E. henningsmoeni the telsons changed from being angular in juveniles to larger and more rounded in adults). All these changes are believed to be a result of the respiratory and reproductive requirements of adults.
Paleoecology
Members of Eurypterus existed for a relatively short time, yet they are the most abundant eurypterids found today. They flourished between the Late Llandovery epoch (around 432 million years ago) to sometime during the Přídolí epoch (418.1 million years ago) of the Silurian period. A span of only around 10 to 14 million years.
During this period, the landmasses were mostly restricted to the southern hemisphere of the Earth, with the supercontinent Gondwana straddling the South Pole. The equator had three continents (Avalonia, Baltica, and Laurentia) which slowly drifted together to form the second supercontinent of Laurussia (also known as Euramerica, not to be confused with Laurasia).
The ancestors of Eurypterus were believed to have originated from Baltica (eastern Laurussia, modern western Eurasia) based on the earliest recorded fossils. During the Silurian, they spread to Laurentia (western Laurussia, modern North America) when the two continents began to collide. They rapidly colonized the continent as invasive species, becoming the most dominant eurypterid in the region. This accounts for why they are the most commonly found genus of eurypterids today. Eurypterus (and other members of Eurypteroidea), however, were unable to cross vast expanses of oceans between the two supercontinents during the Silurian. Their range were thus limited to the coastlines and the large, shallow, and hypersaline inland seas of Laurussia.
They are now only known from fossils from North America, Europe, and northwestern Asia, cratons that were the former components of Laurussia. While three species of Eurypterus were purportedly discovered in China in 1957, the evidence of them belonging to the genus (or if they were even eurypterids at all) is nonexistent. No other traces of Eurypterus in modern continents from Gondwana are currently known.
Eurypterus are very common fossils in their regions of occurrence, millions of specimens are possible in a given area, though access to the rock formations may be difficult. Most fossil eurypterids are the disjointed shed exoskeleton (known as exuviae) of individuals after molting (ecdysis). Some are complete but are most probably exuviae as well. Fossils of the actual remains of eurypterids (i.e. their carcasses) are relatively rare. Fossil eurypterids are often deposited in characteristic windrows, probably a result of wave and wind action.
| Biology and health sciences | Fossil arthropods | Animals |
893280 | https://en.wikipedia.org/wiki/Floriculture | Floriculture | Floriculture is the study of the efficient production of the plants that produce showy, colorful flowers and foliage for human enjoyment in human environments. It is a commercially successful branch of horticulture and agriculture found throughout the world. Efficient production practices have been developed over the years, for the hundreds of plant taxa used in the floral industry, increasing the overall knowledge of whole plant biology. Plant breeding and selection have produced tens of thousands of new genotypes for human use. Jasmine, marigold, chrysanthemum, rose, orchid, and anthurium are flowers of commercial demand.
Overview
Flowers are an important part of human society that are often used at times of joy and sadness, and as a part of everyday life. Flowers and plants may be indoors in a sunny window, as part of the landscape in the front yard or on the patio or deck in the back yard. People have been studying flowers and plants and their interaction with humans and how to produce these flowers and plants so all humans can enjoy them. Floriculture scientists throughout the world to do this work.
Floriculture crops include cut flowers and cut cultivated greens, bedding plants (garden flowers or annuals, and perennials, houseplants (foliage plants and flowering potted plants). These plants are produced in ground beds, flower fields or in containers in a greenhouse. Protected cultivation is often used because these plants have a high value to humans.
Flower crops are grown in simple to highly sophisticated ways. These crops can be grown in soil in farm fields or in field soil in inexpensive high tunnel greenhouses. For years, flowers were grown, seasonally for the specific crop, close to the market in Europe, North America and Asia. However, many crops of the floral industry have moved to a specific climate, typically in the mountains of South America, Africa and China, so certain plants can be grown year around where hand labor is available.
Protected horticulture (greenhouses) has developed simultaneously with the continued changes in the flower crops and markets. Floriculture is a major component of controlled-environment agriculture (CEA). Floriculture crops have a high value to humans, so the cost of an expensive production system - greenhouses, automated environmental control, automated irrigation and fertilization, robotic seed, transplant and container handling, supplemental photosynthetic lighting - is necessary to produce these plants efficiently for the world-wide markets. Some are irrigated manually, but most are irrigated with drip irrigation, boom irrigation or flood floors. Hydroponics can be used for many cut flower crops.
Floriculture value 2022
The global Floriculture market size is estimated to be worth US$50040 million in 2022 and is forecast to be a readjusted size of US$58030 million by 2028 with a compound annual growth rate of 2.5% during the review period.
The total wholesale value of sales across all U.S. floriculture crops totaled US$6.69 billion in 2022 from 8,951 floriculture producers with a production area of 833 million square feet.
Floriculture crops
The horticulture industry involves the cultivation of plants for food, medicine, and aesthetic purposes. It covers a wide range of activities including the growth, distribution, and sale of fruits, vegetables, flowers, trees, and ornamental plants. The sector plays a vital role in agriculture and includes several key branches such as:
1. Floriculture: Focuses on growing and marketing flowers and ornamental plants.
2. Olericulture: The production of vegetables.
3. Pomology: The cultivation of fruit trees and shrubs.
4. Landscape Horticulture: Involves designing, maintaining, and managing landscapes for aesthetic and functional purposes.
Importance of the Horticulture Industry:
Economic Contribution: Horticulture contributes significantly to the global economy by providing employment in both rural and urban areas, and through the sale of fresh produce, ornamental plants, and landscaping services.
Environmental Benefits: It helps improve biodiversity, soil quality, and air purification, as well as reduce urban heat islands through landscaping and green spaces.
Food Security: It enhances food production by supplying fruits and vegetables, which are crucial for balanced diets.
Innovation and Technology: The industry increasingly relies on modern technologies like precision farming, greenhouse systems, and sustainable practices to improve productivity and reduce environmental impact.
The horticulture industry is growing due to the increasing demand for fresh, healthy food, sustainable landscaping, and environmental conservation.
Annual Bedding/Garden Plants
Potted Flowering Plants
Herbaceous Perennial Plants
Foliage Plants, Indoor/Patio Use (Houseplants)
Propagative Floriculture Materials
Cut Flowers
Cut Cultivated Greens
Floriculture advancements
Plant enthusiasts and growers learned significant details about growing certain plants over the years. Chrysanthemums have been cultivated in China for over 3000 years, so growers knew about the plant and how to grow it. Floriculture scientists have simply continued this trend to control the plant's environment to control flowering for the significant dates when humans want flowers for celebrations and gatherings.
Photoperiodism
Chrysanthemum was one of the plants used in experiments that led to the definitions of photoperiod and photoperiodism. Yet, it's likely that Chinese, Korean and Japanese plantsmen had a good understanding based on their years of experience. The occurrence of this physiological response and the reasons for it have been the subject of many experiments at universities and in industry. Poinsettias are another short day plant with importance to flower growers. These and additional experiments and experience have shown that temperature has an impact on the photoperiodic response. Many cut flower and bedding plant species respond to long day or short day treatments for faster flowering. The use of lighting treatments to extend the day and black cloth treatments to shorten the day are important additions to floriculture to increase the efficiency of plant production.
Plant tissue culture, micropropagation
Plant propagation has always been a part of flower and plant gardening. Plant tissue culture began as a way to save orchid embryos as orchid fanciers bred new cultivars. Most horticulture and many botany programs in the world had scientists working on plant propagation through tissue culture techniques from the 1950s to the 1980s. These programs expanded the knowledge base on a wide range of taxa and allowed industry to find the connection to commercial production. Plant tissue culture allowed new, unique phenotypes and genotypes to be propagated in large numbers quickly. Many cultivars of foliage plants are available only from tissue culture. Uniquely, tissue cultured geraniums were heat treated to allow the identification and removal of many viruses, virus-indexed. As viruses were removed, many horticultural characteristics of the many cultivars disappeared; this led plant breeders to leave many viruses in breeding lines for future cultivars. Heat treatment of tissue culture of many taxa has since been used to remove bacteria and virus pathogens in various floriculture crops.
Containers and growing media
Containers of various kinds have been used in the culture of plants for a long time. Field soil or garden soil possibly with an addition of organic matter (compost) was placed in the container or pot and a plant was added followed by regular watering. It required experience and a watchful eye to prevent overwatering.
This success was tied to a relatively deep pot, usually 6–10 inches (15–25 cm) deep or larger. Gravity was sufficient to pull or drain water from the soil so an adequate portion of the soil in the pot was well drained and oxygen would be available to the root system. As US greenhouses began to expand the bedding plant business in the 1950s and 1960s, they needed smaller containers for the logistical aspects of plant spacing and shipping. Vacuum formed plastic trays and packs offered the smaller sizes but composted field soil was easy to overwater in the smaller containers. The first step was to add peat moss and perlite to the field soil in a 1:1:1 ratio. The next step was to use other materials, sphagnum moss peat and vermiculite, in a 1:1 ratio, the Cornell peat-lite mix. In the 1970s, more materials were used for growing media by the companies formed to process and distribute growing media to operations across the country. The physical properties of all the products had to be evaluated on a standard basis to make wise choices with economic decisions the operations were making. As plug (young plant) production, mechanization of seed germination and mechanization of transplanting, began in the 1980s more work was necessary to manage the small volume of growing media in plug trays. Research continues of all aspects of growing media and container design.
The harvest and use of peat for growing media remains an environmental issue in North America and Europe. Alternative and more sustainable materials continue to be added to growing media processing - pine bark, processed pine bark, coco coir, wood fiber, etc. Sustainable solutions for growing media materials remain a high priority for the industry.
Pesticide residues
Pesticide residues remain a significant issue for floriculture crops. Many countries have limited controls on pesticide usage but flower handlers and consumers could be contaminated by the residue.
The impact of certain pesticides, neonics, on bees and other pollinators has become a significant concern. The application of these pesticides on garden flowers during greenhouse production can have a major impact on pollinator populations in a consumer's garden.
Research continues on biological control of greenhouse insect, mite and plant pathogens to reduce pesticide use in floriculture crop production.
Supplemental lighting
Supplemental lighting for flower crops began with photoperiod treatments and interest expanded to determine whether artificial light from electric lamps could substitute for sunlight during winter conditions. Incandescent lamps were not successful, so floriculture had to wait for lighting technology to improve. Advancements with fluorescent lamps and industrial lamps (mercury vapor, high pressure sodium, low pressure sodium, etc.) led to improved plant production for geraniums, roses and other crops. In the following decades, artificial lighting became standard practice in Europe, North America and Japan.
Work was completed to standardize a plant's need for light (radiant energy) from natural and artificial sources. The term daily light integral (DLI) was introduced as a measurement of the optimal amount of radiant energy each plant requires for optimal growth.
The introduction of light emitting diode (LED) lamps offered more opportunities for supplemental lighting. These lamps were more efficient at light production, cooler and allowed the manipulation of light quality from different wavelengths of light compared to other lamps.
Supplemental lighting has been used to optimize production of seedlings, bedding plants, cut flowers and other crops.
Plant nutrition, water quality and irrigation
Flower crops were grown in field soil like all horticultural and agricultural crops. Nutrients important to the flowers were held in the soil matrix and supplemented with additions of organic matter and animal manure. These organic additions were labor-intensive and inconsistent, reducing the ability to optimize flower production. Floriculture moved to growing media and inorganic fertilizer products in the 1950s and 1960s as container production became more important. This move was supported by hydroponic research more than soil science research. The "soil-less" nature of hydroponics was more similar to the "soil-less" nature of growing media.
| Technology | Horticulture | null |
893573 | https://en.wikipedia.org/wiki/Didelphodon | Didelphodon | Didelphodon (from [is] "opossum" tooth "tooth") is a genus of extinct metatherian mammal from the Late Cretaceous of North America.
Description
Although perhaps little larger than a Virginia opossum, with a skull length of and a weight of , the teeth have specialized bladelike cusps and carnassial notches, indicating that the animal was a predator; the jaws are short and massive and bear enormous, bulbous premolar teeth which appear to have been used for crushing. Analyses of a near-complete skull referred to Didelphodon show that it had an unusually high bite force quotient (i.e. bite force relative to body size) among Mesozoic mammals, suggesting a durophagous diet. However, its skull lacks the vaulted forehead of hyenas and other specialized bone-eating durophagous mammals, indicating that its diet was perhaps a mixture of hard foodstuffs (e.g. snails, bones) alongside small vertebrates and carrion; although omnivorous habits were suggested in the past, it appears that it was incapable of processing plant matter, rendering it more likely to be hypercarnivorous or durophagous. Some convergence with the carnassials of other predatory mammal groups has also been noted.
Discovery
Three species of Didelphodon are known: D. vorax, D. padanicus, and D. coyi. The genus is known from the Hell Creek Formation of Montana and the Lance Formation of Wyoming, the Frenchman Formation of Saskatchewan, the Horseshoe Canyon Formation of Alberta, and the Scollard Formation of Alberta, where it is one of the most abundant mammals. It is found solely in late Maastrichtian deposits.and royal tyrrell museuam
Classification
Didelphodon is a stagodontid marsupial related to Eodelphis and Pariadens. The genus appears to descend from the Campanian Eodelphis, and in particular appears to be related to Eodelphis cutleri. Pariadens appears to be more primitive than either Eodelphis or Didelphodon, and is probably sister to their group. Didelphimorphia is an order that was named in 1872 by Gill. Previously, in 1821, Gray named the superfamily Didelphoidea to house the families Alphadontidae, Pediomyidae, Peradectidae, and Stagodontidae, which unites Didelphodon with many other genera. In 2006, a study found that the stagodontids only contained two taxa, Didelphodon and Eodelphis. The previously-included Pariadens was excluded from the group because its type species, P. kirklandi, lacks any of the clade's characteristics; it was reassigned to Marsupialia incertae sedis. Another species, "P." mckennai lacks marsupial features, and is probably a therian. Another historical stagodontid, Boreodon, is a nomen dubium. Finally, the purported stagodontid Delphodon is probably a synonym of Pediomys or Alphadon.
A 2016 phylogenetic analysis found that Didelphodon and other stagodontids were marsupialiforms. Their relationships within the Marsupialiformes are shown below.
Paleobiology
Although it has been argued on the basis of the shape of referred tarsal bones that Didelphodon and other stagodontids were semiaquatic due to having flexible feet, these traits may in fact be evidence of increased rigidity in the foot. Nevertheless, a recently-found and as-of-yet undescribed specimen, located just away from a Triceratops in a riverbed, suggests that Didelphodon may have possessed an otter-like body with a tasmanian devil-like skull. A study that is being prepared by Kraig Derstler, Greg Wilson, Robert Bakker, Ray Vodden and Mike Triebold will describe this new specimen, housed in the Rocky Mountain Dinosaur Resource Center. A study on Mesozoic mammal locomotion demonstrates that Didelphodon groups with semi-aquatic species.
The evolution of Didelphodon and other large stagodontids (as well as large deltatheroidans like Nanocuris) occurs after the local extinction of eutriconodont mammals, suggesting passive or direct ecological replacement. Given that all insectivorous and carnivorous mammal groups suffered heavy losses during the mid-Cretaceous, it seems likely these metatherians simply occupied niches left after the extinction of most eutriconodonts.
| Biology and health sciences | Marsupials | Animals |
893602 | https://en.wikipedia.org/wiki/Ambulocetus | Ambulocetus | Ambulocetus (Latin ambulare "to walk" + cetus "whale") is a genus of early amphibious cetacean from the Kuldana Formation in Pakistan, roughly 48 or 47 million years ago during the Early Eocene (Lutetian). It contains one species, Ambulocetus natans (Latin natans "swimming"), known solely from a near-complete skeleton. Ambulocetus is among the best-studied of Eocene cetaceans, and serves as an instrumental find in the study of cetacean evolution and their transition from land to sea, as it was the first cetacean discovered to preserve a suite of adaptations consistent with an amphibious lifestyle. Ambulocetus is classified in the group Archaeoceti—the ancient forerunners of modern cetaceans whose members span the transition from land to sea—and in the family Ambulocetidae, which includes Himalayacetus and Gandakasia (also from the Eocene of the Indian subcontinent).
Ambulocetus had a narrow, streamlined body, and a long, broad snout, with eyes positioned at the very top of its head. Because of these features, it is hypothesised to have behaved much like a crocodile, waiting near the water's surface to ambush large mammals, using its powerful jaws to clamp onto and drown or thrash prey. Additionally, its ears possessed similar traits to modern cetaceans, which are specialised for hearing and detecting certain frequencies underwater, although it is unclear if Ambulocetus also used these specialised ears for hearing underwater. They may have instead been utilised for bone conduction on land, or perhaps served no function for early cetaceans.
It is thought to have swum much like a modern river otter, tucking in its forelimbs while alternating its hind limbs for propulsion, as well as undulating the torso and tail. It may have had webbed feet, and unlike its modern relatives, lacked a tail fluke. On land, Ambulocetus may have walked much like a sea lion.
Ambulocetus inhabited the Indian subcontinent during the Eocene. The area had a hot climate with tropical rainforests and coastal mangroves, and Ambulocetus may have predominantly inhabited brackish areas such as river mouths. It lived alongside requiem sharks, catfish and various other fishes, turtles, crocodiles, the amphibious hoofed mammal Anthracobune, and the fellow cetaceans Gandakasia, Attockicetus, Nalacetus, and Pakicetus.
Taxonomy
Discovery
In December 1991, Pakistani palaeontologist Mohammad Arif and Dutch–American palaeontologist Hans Thewissen were jointly funded by Howard University and the Geological Survey of Pakistan to recover land mammal fossils in the Kala Chitta Hills of Punjab, Pakistan. On 3 January 1992, they recovered a small, thick rib fragment. Later in the field season, while surveying the upper Kuldana Formation, Thewissen discovered a femur (thigh bone) and proximal portion of the tibia (upper portion of the shin) which clearly belonged to a mammal. An hour later, Arif discovered the rest of the skeleton, and the two began excavation the next day. At first, Thewissen speculated the fossils belonged to an anthracobunid (a large semi-aquatic mammal), until he found the teeth near the end of the field season, which were characteristically cetacean (living cetaceans are whales, dolphins, and porpoises). Thewissen, at the time, could not afford to excavate and store everything, so he took the skull with him to the United States, while Arif kept the rest in two crates which used to hold oranges. In October 1992, Thewissen presented his research of the skull to a vertebrate palaeontology convention in Toronto, Canada. The next year, American palaeontologist Philip D. Gingerich paid for the rest of the skeleton to be shipped to the United States. In 1994, the formal description of the remains was published by Thewissen, mammal palaeontologist Sayed Taseer Hussain, and Arif. They identified the remains as clearly belonging to an amphibious cetacean, and so they named it Ambulocetus natans. The genus name comes from Latin ambulare "to walk" and cetus "whale", and the species name natans "swimming".
The Kuldana Formation is constrained to sometime during the Lutetian stage of the Early Eocene, and the remains may date to 48–47 million years ago. The holotype specimen, HGSP 18507, is a partial skeleton initially discovered preserving an incomplete skull (missing the snout), some elements of the vertebral column and ribs, as well as portions of the fore- and hind-limb. Other specimens initially found were HGSP 18473 (a second premolar), HGSP 18497 (a third premolar), HGSP 18472 (a tail vertebra), and HGSP 18476 (lower portion of a femur). The holotype was found in a silt and mudstone bed over a area. Further excavation recovered most of the holotype's skeleton—most notably the hip, sacrum, and most of ribcage and thoracolumbar series (the spine excluding the neck, sacrum, and tail). These left the holotype about 80% complete by 2002, making it the most completely known cetacean from the time period. In 2009, some more elements of the holotype's jawbone were identified from a then-recently prepared matrix block.
Though it was known that cetaceans descended from land mammals before the discovery of Ambulocetus, the only evidence of this in the fossil record was the 52-million-year-old (fully terrestrial) Pakicetus and the Paleocene mesonychians (as there was a hypothesised link between cetaceans and mesonychians). The limbs of more aquatic Eocene cetaceans did not preserve very well. Ambulocetus demonstrated that cetaceans swam by flexing the spine up and down (undulation) before they had evolved the tail fluke, forelimb propulsion evolved relatively late, and that cetaceans went through an otter-like phase with spinal undulation and hindlimb propulsion. These had already been hypothesised to have occurred in the earliest aquatic cetaceans, but were impossible to test without more complete remains. The describers noted that, "Ambulocetus represents a critical intermediate between land mammals and marine cetaceans."
Classification
Modern cetaceans (Neoceti) are grouped into either the parvorders Mysticeti (baleen whales) or Odontoceti (toothed whales). Neoceti are descended from the ancient Archaeoceti, whose members span the transition from terrestrial to fully aquatic. Archaeoceti are thus paraphyletic (it is a non-natural group which does not comprise both a common ancestor and all of its descendants). Ambulocetus was an archaeocete. By the time Ambulocetus was discovered, archaeocetes were classified into the families Protocetidae (which included what are now the terrestrial Pakicetidae, and the rest were amphibious), Remingtonocetidae (amphibious), Basilosauridae (aquatic), and Dorudontidae (aquatic, now a subfamily of Basilosauridae). The earliest cetaceans were thought to be the mesonychians, proposed before any firm early cetacean fossils were identified. In the original description, Ambulocetus was preliminarily placed into Protocetidae, until the further description of the holotype prompted Thewissen and colleagues to move it into its own family Ambulocetidae in 1996. At the same time, they also erected the family Pakicetidae. They also proposed that some members of Pakicetidae, Protocetidae, and Ambulocetidae were the other two archaeocete families' ancestors. They suggested that mesonychians gave rise to pakicetids, which gave rise to ambulocetids, which gave rise to both protocetids and remingtonocetids.
Though middle-to-late-Eocene archaeocetes are also known from North America, Europe, and Africa, most of these are found only on the Indian subcontinent. Therefore, it is thought cetaceans originally evolved in that region. Based on molecular data, cetaceans are most closely allied with hippos (Whippomorpha), and they split approximately 54.9 million years ago. They are all placed in the order Cetartiodactyla alongside terrestrial even-toed ungulates (hoofed mammals). This puts mesonychians as a distant relative of cetaceans rather than an ancestor, and their somewhat similar morphology was possibly a result of convergent evolution. The oldest known cetacean is the ambulocetid Himalayacetus identified in 1998 and dated to 52.5 million years ago (predating the terrestrial pakicetids), though the exact dating of Himalayacetus and Pakicetus is debated. Ambulocetidae also includes Gandakasia. Himalayacetus and Gandakasia are known only from partial jaw fragments. Ambulocetidae are endemic to the Indian subcontinent, and span the early to middle Eocene.
Description
Size
Upon description, Thewissen and colleagues suggested the holotype specimen may have weighed the same as a male South American sea lion — about — based on the size of the vertebrae, ribs, and limbs. They also estimated a length of roughly . For comparison, the holotype of Pakicetus attocki may have been long. In 1996, they estimated weight of Ambulocetus, using the cross-sections of the long bones, as . Alternatively, they estimated about by using the length of the second upper and lower molars compared to trends between this length and ungulate body mass. They obtained the same result comparing the skull size to those of similarly sized carnivores. In 1998, based on vertebral size, Gingerich estimated a body mass of , similar to modern cetaceans. In 2013, Thewissen suggested that this may be an unreliable mass determinant as the vertebrae are unusually robust in Ambulocetus.
Skull
Like other archaeocetes which preserve this element, the base of the skull has an undulating contour, probably related to the shape of the nasal canal (and its passage to the throat) and the narrow infraorbital region (the area below the eyes). The base of the skull is wide compared to other archaeocetes, more like that of modern cetaceans. The narrow infraorbital space, made of primarily the pterygoid processes, also occurs in Remingtonocetus and Pakicetus. The pterygoids connect as far back as the middle ear, much farther than other archaeocetes including the more ancient Pakicetus. Most modern cetaceans have a falcate (sickle-shaped) process which juts out prominently halfway between the hypoglossal canal and the ear; Ambulocetus has a similar process continuous of the pterygoid, but it runs alongside and behind the hypoglossal canal. Like many other archaeocetes, the pterygoids, sphenoids, and palatines form a wall lining the bottom of the nasal canal, which causes the palate to extend all the way to the ear. Like other cetaceans, Ambulocetus lacks the postglenoid foramen, which usually is one of the main passageways for veins into the skull in placental mammals. The ectotympanic bone which supports the eardrum is similar to that of Pakicetus, about as long as wide, whereas later archaeocetes have more elongate ectotympanics. The ectotympanics of all archaeocetes, nonetheless, are much different than those of terrestrial mammals. The ectotympanics of all cetaceans, including Ambulocetus, possess an involucrum (thickened lump of bone) at the medial lip. Unlike Pakicetus, but like later archaeocetes, the tympanic made close contact with the jaw. Like later archaeocetes, Ambulocetus seems to have possessed an air sinus in the pterygoids. It may have also had paranasal sinuses. The parietal bones on the braincase sides are more perpendicular than in Remingtonocetus, which makes the cheeks appear less flared. Like Remingtonocetus, Ambulocetus appears to have had a small brain.
The snout was quite broad, but the end of the holotype's snout is missing, so it is unclear how long it would have been. The snouts of Basilosaurus and Rodhocetus are short and make up about half the skull's length. Remingtonocetid snouts are quite narrow, which was clearly not the case for Ambulocetus. The mandibular symphysis of most mammals is restricted to the midline of the jaw, but extends much farther in archaeocetes; in Ambulocetus, it reaches the back end of the first premolar. Snout robustness and symphysis length suggest reinforcement of the jaw to withstand a strong bite force. Similarly, the strongest biting muscle in Ambulocetus seems to have been the temporalis muscle involved in biting down. Like other cetaceans, there are embrasure pits (a depression between the teeth), preserving the tooth positions for the fourth premolar, the first molar, and the third molar. Unlike later archaeocetes, the molars' roots do not extend to the cheek bones, and the third molar is not as nosewards as in remingtonocetids. The coronoid process of the mandible (where the lower jaw connects with the skull) in Ambulocetus is steep. In contrast, it is low and slopes gently down in basilosaurids and later cetaceans. The mandibular foramen opens below the coronoid process, and is around midway between terrestrial mammals and toothed whales in size. Like other cetaceans, the body of the hyoid bone (the basihyoid bone) is about as long as wide. Unlike other archaeocetes, the eyes are quite large and are placed near the top of the head facing upwards.
Unlike modern toothed whales which only have one kind of tooth (homodont), archaeocetes are heterodont. Judging by tooth root size, the lower canine was larger than the incisors. The teeth are more robust than those of Rodhocetus and Basilosaurus. The premolars were double rooted, whereas most archaeocetes have single-rooted first premolars. The enamel of the lower premolars is crenulated (has scalloped edges). The fourth premolar is a high triangular shape. Like other ancient cetaceans, and most pronouncedly in ambulocetids, the lower molars are shorter than the back premolars. The lower premolars are larger than those of Pakicetus and are separated by wider gaps (diastemata). The molars had distinct trigonid and talonid cusps (these cusps are lost in basilosaurids), and the upper molars were trituberculate like ancient archaeocetes and ancient placental mammals, meaning they had a large protocone, distinct paracone and metacone, and no accessory cusps. Later archaeocetes developed accessory cusps.
Ribs and vertebrae
The holotype preserved seven neck vertebrae, which are rather long at . The 16 preserved thoracic vertebrae have thick spinous and transverse processes (which jut upwards and obliquely from the centrum, the vertebral body), with deep depressions on both sides at the tail-end of each centrum which may have supported strong longissimus muscles which flex the spine. The thoracic vertebrae become longer and wider tailwards and are tallest mid-series. In front-view (anterior aspect), the centra go from heart-shaped to kidney-shaped by T8 (the eighth thoracic vertebra). The pedicals (between the centrum and a transverse process) feature deep grooves. The spinous processes project tailwards from T1–T9, straight up at T10, headwards from T11 to T12, and the rest project straight up. The spinous processes progressively increase in length and width from T11–T16. T10 seems to have been at the level of the thoracic diaphragm. T1–T12 and T14 have capitular facets on the top margin of both the frontward and tailward side to join with the ribs. T15 and T16 have capitular facets on the headward side and lack transverse processes. T11–T15 have accessory anapophyses which jut straight up from the top border between the centrum and the transverse processes; and in T16, these are small, originate near the pedicles, and project tailwards. The width between articular processes (two masses of bone which jut out of each centrum to connect with the next centrum) continually increases through the thoracolumbar series. In life, it is possible it had up to 17 thoracic vertebrae.
The holotype preserves 26 ribs, though it is thought to have had 32 total in life. The cortical bone (the outermost layer) is thickest at the neck of the rib (between the joint and the costal cartilage), at max , and was filled with spongy bone. That is, unlike many other aquatic mammals, the ribs did not exhibit osteosclerosis. They did exhibit pachyostosis, and were made thicker and heavier with additional layers of lamellar bone. The ribs' shape indicates Ambulocetus had a narrow and heart-shaped thorax looking at it head-on. Ribs are thickest at the T8–T10 level. Ribs are broadest at the sternum, which suggests strong sternocostal joints. The ribs have a slight S-curve in side view, with the rib heads angled headwards, and the sternocostal joints angled tailwards. The holotype preserves a central and a tailward sternum bone which are both exceedingly thick, about on the outer margins and decreasing towards the centre. The central sternum bone is longer and wider than the tailward one.
The eight preserved lumbar vertebrae at the lower back are much longer than the thoracic, and the centra and transverse processes, from L1–L7, continually increase in length and height. The short transverse processes on L8 are probably due to its proximity to the ilium on the hip. The undersides are concave. The spinous processes are long and tall, and project headward from L1–L5, and straight-up from L6–L8. The spinous processes are bulbous on the tailward side to support epaxial muscles. The vertebral laminae are excavated headward to support the interspinous ligaments which connect the spinous processes. The vertebrae are about as robust as those of modern pinniped such as sow leopard seals and bull walruses. The postzygapophyses (the surface where the vertebrae join with each other) is flat rather than revolute, which would have made the series more flexible than that of terrestrial relatives.
For the four preserved sacral vertebrae (at the sacrum, between the pelvic bones), the transverse processes of S1 are smaller than those of L8. There is a robust sacroiliac joint with the hip. For the spinous processes, those of S1–S3 are fused. Metapophyses jut straight up from each lamina near the joint, progressively getting smaller with each vertebra.
Only five of the tail (caudal) vertebrae are preserved: a possible C1 or C2, a possible C3, a possible C4, a possible C7, and a possible C8. The more headward tail vertebrae have thick transverse processes, whereas those of the middle tail vertebrae are longer than broad. The C3 has a narrow spinous process and is mostly columnar, but the tailward side is broader. The C4 is more columnar. The C7 and C8 are columnar and taper off tailward, and the neural canal where the central nervous system runs through is still present. In life, Ambulocetus possibly had upwards of 20 tail vertebrae like some mesonychians. If correct, then Ambulocetus would have had a lot fewer and a lot longer tail vertebrae than modern cetaceans.
Limbs and girdles
Unlike modern cetaceans, Ambulocetus had functional legs which could support the animal's bodyweight on land. The holotype has a robust radius and ulna (the forearm bones). The forearm measures in length. The head of the radius is somewhat triangular, which probably means the forearm was locked in a semi-pronated position (the palms were orientated towards the ground). The olecranon, which formed part of the elbow joint, makes up about a third of the ulna's length and is inclined tailwards, which would have allowed the triceps to more forcefully flex the elbow. The wrist bones indicate a strong flexor carpi ulnaris muscle for wrist flexion. The hand had five widely spaced digits. The first metacarpal (which is in the thumb) is long, the second , the third , the fourth , and the fifth . Like modern beaked whales, the thumb is short and slender.
The ilium of the hip of Ambulocetus, like remingtonocetids, features deep depressions to support the rectus femoris and the gluteal muscles. Unlike terrestrial mammals and protocetids, the ischium is expanded dorsolaterally (from left to right, and upwards), which would have increased lever arm for thigh and leg retractor muscles when extended, such as while swimming. This would have also increased the surface area of the gemelli muscles (hip rotators which stabilise the hip) and the tail muscles. The widening of the ischium may have also given Ambulocetus a more streamlined and hydrodynamic body. Ambulocetus had a pubic symphysis connecting the two pubic bones at the base of the pelvis together, which indicates the animal could support its own weight on land. The modern cetacean pubis bone lacks this and only functions to anchor abdominal and urogenital muscles.
The leg proportions of Ambulocetus are similar to otters and seals, and American mammalogist Alfred Brazier Howell predicted similar proportions for a transitional cetacean in 1930. The femur measures , a length similar to the presumably cursorial (capable of running) mesonychian Pachyaena. Archaeocete femora are generally much shorter. The femoral head is spherical and, at maximum, has a width of , similar to Indocetus but much larger than mesonychians and Rodhocetus. The trochanteric fossa, supporting the lateral rotator group at the hip, is quite deep, but other than this, the femur does not seem to have supported particularly strong extensor or flexor muscles. The femoral condyles of Ambulocetus are quite long compared to those of other archaeocetes and mesonychians, suggesting the knee was capable of hyperflexion (bending). The tibia is overall similar to those of mesonychians. The feet are huge, probably longer than the femur. The toes are also relatively long, with the fourth digit measuring in length. The fifth digit is slightly shorter and much less robust than the fourth. The phalanges of the toes are short, and end with a convex hoof. Like seals, the phalanges of both the hands and feet are flattened, which may have streamlined them to allow for webbed feet.
Palaeobiology
Diet
The robustness of the cheek teeth, as well as the cusp arrangement, suggests they were involved in crushing, and the fact that both the premolars and molars were involved in crushing indicates Ambulocetus required a large area for crushing, such as when biting into large prey items. Similarly, the broad and powerful snout makes it unlikely it was pursuing small, quick prey items (which would have required a narrow snout like dolphins or gharials). The snout was also long, which may have precluded the ability to crush bone because it would have had reduced structural integrity at the tip. The anatomy of the cheek teeth resembles those of Mesozoic marine reptiles which fed on armoured fish, large fish, reptiles, and ammonites, and the teeth may have been used to grip onto prey firmly. Therefore, Thewissen suggested Ambulocetus was most likely an ambush predator, the jaw adapted to handle struggling prey. The unusually deep pterygoids potentially functioned to dissipate force while the prey was struggling.
The eyes of Ambulocetus were placed on the top of the head, similar to crocodiles and other animals that prefer to keep most of their body submerged with the eyes peeking out of the water. The nasal canal has bony walls extending into the throat, much like in crocodiles where they keep the nasal airways open while the animal is killing prey either by drowning it or thrashing it around. Pieces of prey are subsequently torn off by forceful, thrashing head and body motions, the feet anchoring the crocodile in place. Thewissen believed Ambulocetus used a similar feeding tactic, though Ambulocetus was probably capable of chewing, unlike crocodiles. Ambulocetus may have attacked large mammals which approached the water's edge, and semi-aquatic mammals including early (possibly herbivorous) sirenians (now manatees and the dugong) and the probably amphibious anthracobunids. These two seem to have been rather common on the coasts of the Indian subcontinent, which could mean they were regular prey items. Since Ambulocetus was found in marine deposits (where animals would not come to drink), it is possible it hunted in river deltas which were recorded in the Kuldana Formation. Ambulocetus probably went after fish and reptiles when given the opportunity, though it may not have had the agility to commonly catch them.
Locomotion
Thewissen hypothesised that Ambulocetus was a drag-powered swimmer, and used its huge feet as its primary propulsion mechanism, much like modern river otters including the giant otter, and species in the genera Lontra and Lutra. Based on the length of the known tail vertebrae, Ambulocetus may have had an inflexible tail, which would have made the tail an inefficient primary propulsion mechanism due to poorer lever arm (modern cetaceans have relatively short tail vertebrae). Ambulocetus therefore likely did not have a tail fluke. Nonetheless, drag powered swimmers still have powerful tails for producing lift, and the tails of river otters are 125% the size of the thoracolumbar series. So, using river otters as a model, Ambulocetus was a pelvic paddler—swimming with alternating beats of the hindlimbs (without engaging the forelimbs)—and also undulated (moved up and down) its tail while swimming. Like the sea otter, pelvic paddling may have been done at the surface to move at slow or moderate speeds. At higher speeds fully submerged, undulation of the spine would have become more prominent, though the feet still would have acted as the primary propulsion mechanism.
Based on the pelvis and robust forelimbs, Thewissen believed Ambulocetus was capable of venturing onto land, and was more efficient at doing this than remingtonocetids and protocetids (it is unclear if the latter two were capable of bearing weight on the limbs). Ambulocetus possibly used a sprawling gait on land, similar to modern sea lions. In 2016, Japanese biologists Konami Ando and Shin-ichi Fujiwara performed a statistical test of ribcage strength among terrestrial, semi-aquatic, and fully aquatic mammals, and found that Ambulocetus clustered with fully aquatic mammals, because they assigned a very high rib density on par with fully aquatic sirenians which use their heavy, osteosclerotic ribs as ballast. They then concluded Ambulocetus could not walk on land, but cautioned the study was limited by a lack of information on the exact density of the bone, the location of the centre of mass, and the reliance of false ribs for thoracic support.
Hearing
Modern cetaceans have highly specialised ear bones to hear underwater as well as to detect certain frequency ranges. Unlike most other mammals, cetacean ear bones are comparatively thick, and so preserve more reliably in the fossil record. Modern cetaceans have air sinuses surrounding the ear bones (peritympanic sinuses), which acoustically isolate the ear by reflecting sound moving through the head and interrupting both bony and fleshy connections of the ear to the skull. Like later archaeocetes, Ambulocetus had at least one such sinus between the tympanic bone and the skull base. The evolution of these sinuses also seems to have caused some restructuring of the skull base due to the development of bony walls surrounding the sinuses. The ectotympanic of all cetaceans, including Pakicetus and Ambulocetus, has a bony growth (involucrum) on the medial lip speculated to aid in the detection of low-frequency sounds. All cetaceans also have a vertical crest ("sigmoid process") right in front of the ear canal, which is speculated to be related to the increasing size of the malleus bone in the middle ear.
As for the outer ear, terrestrial mammals channel sound in via an ear canal, but those of modern cetaceans are either narrowed or completely plugged, the sound being picked up (at least for toothed whales) by a fat pad in the lower jaw running to the ectotympanic bone. The mandibular foramen size can determine the size of the fat pad, and that of Ambulocetus is larger than that of Pakicetus and terrestrial mammals, but is smaller than later archaeocetes and toothed whales. Nonetheless, a lot of the change to the external auditory apparatus occurred between Pakicetus and Ambulocetus. These early archaeocetes may have developed such an external ear to either: better hear underwater; facilitate bone conduction of vibrations on dry land as some low-lying terrestrial creatures do (namely turtles and subterranean mole rats); or it was non-functional, and the malleus and jawbone (which are connected in the embryo stage of mammals) happened to stop separating.
Palaeoecology
During the Eocene, the Indian subcontinent was an island just beginning its collision with Asia which would eventually lead to the uprising of the Himalayas. The Eocene had a greenhouse climate (no permanent ice sheets at the poles) as opposed to the icehouse climate of today, so, in general, areas were much warmer. The abundance of Eocene brown coal deposits preserving tropical biota on the Indian subcontinent indicates the proliferation of tropical rainforests in a hot climate. Mangroves seem to have commonly grown along the subcontinent's western margin in the Early Eocene but decreased nearing the Middle Eocene Climatic Optimum (a warming trend). The waters off the western coast seem to have featured upwelling and low oxygen.
The holotype was identified in the upper level of the Kuldana Formation at Locality 9209, which features green mud and silt as well as a bed of marine shells, including marine snails (such as Turritella) and bivalves. It was likely a coastal area. A redbed underlies this layer, followed by grey, green, and purple freshwater mud, silts, sandstones, and limestone. These beds alternate with showing marine deposits. The holotype was found in green mud. Near Locality 9209, the formation begins with of grey and green mud, silt, and sandstone, containing two bivalve beds. The first often stretches only one shell, whereas the second stretches down, and the formation terminates with a bed before transitioning to the younger Kohat Formation. The holotype was found a few decimeters above the second bed.
The area may have formed in a shallow sea off the shores of a coastal swamp or forest. The only other vertebrate remains found at the 9209 locality was a (now lost) reptile scute. Other localities of the upper level of the formation have yielded remains of requiem sharks, the fish Stephanodus, catfish, turtles, crocodiles, and the anthracobunid Anthracobune pinfoldi. Other archaeocetes from the formation are: the ambulocetid Gandakasia, the remingtonocetid Attockicetus, and the pakicetids Nalacetus, Pakicetus calcis, and P. chittas. Stable carbon and oxygen isotope analysis indicates Ambulocetus inhabited brackish waters (part fresh and part salt water), possibly at a river mouth.
| Biology and health sciences | Cetaceans | Animals |
894250 | https://en.wikipedia.org/wiki/Orthopedic%20cast | Orthopedic cast | Orthopedic casts or just casts are a form of medical treatment used to immobilize and support bones and soft tissues during the healing process after fractures, surgeries, or severe injuries. By restricting movement, casts provide stability to the affected area, enabling proper alignment and healing of bones, ligaments, and tendons. They are commonly applied to the limbs but can also be used for the trunk, neck, or other parts of the body in specific cases. Orthopedic casts come in various types and designs, tailored to the nature and severity of the injury, as well as the patient's needs. Advances in medical techniques have made casts more comfortable, effective, and versatile, allowing for both weight-bearing and non-weight-bearing options.
Upper extremity casts
Upper extremity casts are frequently utilized to immobilize the arm, wrist, or hand for the treatment of fractures, soft tissue injuries, or during post-surgical recovery. They offer stabilization and support, aiding in proper healing while minimizing the risk of further injury. Common types include long arm casts, short arm casts, and specialized versions such as thumb spica casts.
Long arm casts
A long arm cast extends from the upper arm to the wrist or hand, immobilizing the elbow joint in addition to the forearm. It is typically used for injuries requiring stabilization across multiple joints, such as forearm fractures, certain elbow injuries, and complex soft-tissue damage. It is usually insured that the elbow remains immobilized in a slightly flexed position, usually around 90 degrees, to promote healing while maintaining comfort. Patients with long arm casts often require close monitoring for swelling and circulation issues, given the cast’s extensive coverage.
Short arm casts
A short arm cast is designed to immobilize the wrist and part of the forearm, extending from below the elbow to the hand, often leaving the fingers free for limited mobility. It is used to treat less severe injuries, such as wrist fractures, sprains, or carpal bone issues. The cast restricts wrist movement while allowing elbow mobility, providing a balance between immobilization and functionality. In some cases, a thumb spica variant is added to include the thumb in immobilization, such as for scaphoid fractures or severe thumb sprains. Proper fit and careful alignment are critical to ensure effective healing and prevent complications.
Lower extremity
Leg casts are designed to immobilize the lower limb, facilitating the healing process for fractures, ligament injuries, or post-surgical repairs. They provide stability to the affected area, helping to alleviate pain and prevent additional damage. The design of leg casts can vary to cater to specific injuries, from simple foot fractures to more complicated multi-joint issues. The most common types of lower extremity casts include, long leg casts, Short leg casts. Different varieties exist between the two main types.
Long leg casts
A long leg cast extends from the upper thigh to the toes, immobilizing the knee joint as well as the lower leg and ankle. It is typically used for injuries requiring stabilization across multiple joints, such as tibial or fibular fractures, severe knee injuries, or post-surgical recovery. It is ensured that the knee remains immobilized in a slightly flexed position, typically around 20-35 degrees, to promote healing while maintaining comfort. Patients with long leg casts often require close monitoring for swelling, circulation issues, and mobility challenges due to the cast’s extensive coverage.
Short leg casts
The short leg cast is designed to immobilize the lower leg and ankle, extending from just below the knee to the toes. It is used to treat less severe injuries, such as ankle fractures, foot fractures, or severe sprains. The cast restricts ankle movement while allowing knee mobility. In some cases, a toe plate is added to a short leg cast to provide additional protection for toe injuries or fractures. The toe plate is an extension of the cast that covers the toes, shielding them from external forces and reducing the risk of further injury during recovery. It also helps maintain proper alignment of the toes, preventing displacement of fractured bones or soft tissue damage. Toe plates are particularly useful for injuries where direct impact or accidental movement could hinder healing, such as complex fractures or severe soft tissue injuries in the toes.
Ambulation
When a patient is advised not to put weight on an injured limb, mobility aids like crutches, walkers, or wheelchairs can be used to help with movement during the recovery process. These aids protect the injured area while allowing the patient to move around safely. For those who cannot use crutches due to balance or strength challenges, a wheelchair or knee scooter might be suggested as a more stable mobility option. In certain situations, partial weight-bearing may be permitted, and specialized footwear can be fitted over the cast for added support. For leg casts that allow weight-bearing, the under-sole is usually reinforced to evenly distribute pressure and minimize strain on the injury. Walking casts, as they are called, come with a hard, flat sole to aid in walking while ensuring proper alignment and stability. Other alternative for ambulation with an injured leg include using an
Specialized and rarely used casts
Cylinder
Cylinder casts are orthopedic devices used to immobilize the arm or leg while leaving the surrounding joints free, providing focused stabilization to specific regions. In the arm, a cylinder cast typically extends from the upper arm to just above the wrist, stabilizing injuries like isolated humeral fractures or post-surgical repairs that do not require elbow immobilization. For the leg, the cast extends from the thigh to just above the ankle, often used to manage patellar fractures, some types of tibial plateau injuries, or post-operative care following knee surgeries. The application involves precise alignment to maintain proper positioning of the affected area while ensuring adjacent joints remain mobile, allowing for some functional movement and reducing stiffness during recovery.
Body
Body casts, also known as full-body casts are devices designed to immobilize the trunk of the body, sometimes extending to the neck, head, or extremities. They are less commonly used today due to advances in less restrictive bracing systems and surgical techniques but remain crucial in specific cases where maximum immobilization is essential. Body casts are often used in pediatric patients, particularly young children, who may struggle to comply with wearing a back brace consistently. They are also employed after radical surgeries to repair spinal injuries, congenital deformities, or significant trauma to the spine, pelvis, or upper thigh.
A common variant, the body jacket, encases the trunk and includes shoulder straps to provide added stabilization, particularly for injuries involving the thoracic or lumbar spine. These casts are meticulously shaped to maintain spinal alignment and prevent movement that could disrupt healing. Despite their efficacy, body casts can be extremely uncomfortable due to their restrictive nature and the challenges they pose for hygiene and mobility. In some cases, openings or windows are incorporated into the cast to allow access for wound care or medical monitoring.
EDF cast
An EDF (elongation, derotation, flexion) cast is a specialized orthopedic device used in the treatment of Infantile Idiopathic Scoliosis. This method of correction was pioneered by UK scoliosis specialist Min Mehta and is a non-surgical approach designed to guide spinal growth and alignment during a critical developmental period. Scoliosis is a three-dimensional spinal deformity requiring correction in all planes—coronal, sagittal, and axial—and EDF casting addresses these complex needs. By employing traction, the EDF method elongates the spine, derotates the vertebrae and pelvis, and improves lordosis and overall body alignment, significantly enhancing the patient’s posture and physical function.
The EDF casting technique is distinct from Risser casting in its design and application. EDF casts are tailored to each child’s anatomy, with configurations either over or under the shoulder, depending on the curve pattern and severity. A key feature is the large mushroom-shaped opening on the front, which facilitates proper chest expansion and breathing. On the back, a small cutout is strategically placed on the concave side of the curve, stopping at the midline. This cutout has been shown to improve rotational correction and enhance spinal alignment compared to casts without it. The combination of elongation, derotation, and flexion in this casting method offers an effective early intervention to correct scoliosis and guide proper spinal development.
Spica cast
A spica cast encases the trunk of the body and one or more limbs, providing immobilization for injuries or conditions requiring stabilization across multiple joints. Spica casts can be used for both upper and lower extremities. For instance, a shoulder spica covers the trunk of the body and one arm, typically extending to the wrist or hand. These casts were once common for severe shoulder injuries but are rarely used today, as specialized splints and slings have largely replaced them, promoting early mobility to prevent joint stiffness during recovery.
A hip spica cast, by contrast, is used to immobilize the trunk and one or more legs. Variants include the single hip spica, which covers the trunk and one leg down to the ankle or foot; the double hip spica, which covers the trunk and both legs; and the one-and-a-half hip spica, which encases one leg fully and the other only to above the knee. The extent of trunk coverage depends on the specific injury or condition and the surgeon’s preference, ranging from the navel for spinal mobility to as high as the rib cage or armpits in rare cases. Hip spicas are commonly used to maintain reduction of femoral fractures, treat congenital hip dislocations in infants, and stabilize the hips and pelvis after surgery.
In some cases, a spica cast may not fully encase the legs, extending only to above the knee. These casts, known as pantaloon casts, are occasionally used to immobilize the lumbar spine or pelvis. When applied for such injuries, the trunk portion of the cast typically extends to the armpits to ensure adequate stabilization.
Hygiene
Maintaining proper hygiene while wearing a cast is crucial to ensure patient comfort, prevent skin irritation, and reduce the risk of infection. Since casts are often made of non-breathable materials and remain in place for weeks, they can create an environment prone to moisture buildup, which can lead to odors, skin irritation, or fungal growth. Patients are advised to keep the cast completely dry, as moisture can weaken the cast material and compromise its integrity. Waterproof covers or plastic bags secured with elastic can be used during bathing to protect the cast, but immersing the cast in water should always be avoided unless it is specifically designed to be waterproof.
To maintain hygiene around the cast, patients should clean and moisturize the exposed skin near the cast edges, being cautious not to let any liquids seep inside. A damp cloth with mild soap can be used for cleaning, followed by gentle drying. Avoid inserting objects, such as sticks or sharp items, into the cast to alleviate itching, as this can cause skin abrasions or damage the cast lining. If itching becomes unbearable or if there is persistent discomfort, a healthcare provider should be consulted rather than attempting to adjust the cast.
For long-term casts, regular inspections by a medical professional are recommended to ensure the skin underneath remains healthy and the cast fits properly. Unpleasant odors, excessive itching, or discharge from the cast are potential signs of an infection or skin breakdown, requiring immediate medical attention. Maintaining proper cast hygiene not only contributes to physical comfort but also supports a safe and successful healing process.
Casting materials
Casts typically come in two main types of material, fiberglass, and plaster, though it is less common. Plaster casts have several limitations, including weight, which restricts movement, and skin complications such as dryness, itching, rashes, and infections, particularly in hot weather. Plaster can also break down if exposed to moisture. The cast removal process, which involves a noisy oscillating saw, can cause distress, especially in children, though it is generally painless. Due to these drawbacks, fiberglass casts were developed in the 1970s, offering a lighter, more durable, and water-resistant alternative, though they still have limitations in terms of skin irritation and moisture management.
Plaster casts
Plaster casts consist of a cotton bandage that has been combined with plaster of paris, which hardens after it has been made wet. Plaster of Paris is calcined gypsum (roasted gypsum), ground to a fine powder by milling. When water is added, the more soluble form of calcium sulfate returns to the relatively insoluble form, and heat is produced.
2 (CaSO4·½ H2O) + 3 H2O → 2 (CaSO4.2H2O) + Heat
The setting of unmodified plaster starts about 10 minutes after mixing and is complete in about 45 minutes; however, the cast is not fully dry for 72 hours.
Fiberglass casts
Bandages of synthetic materials are also used—often knitted fiberglass bandages impregnated with polyurethane, sometimes bandages of thermoplastic. These are lighter and dry much faster than plaster bandages. However, plaster can be more easily moulded to make a snug and therefore more comfortable fit. In addition, plaster is much smoother and does not snag clothing or abrade the skin.
Cast liner
Traditional cast liners are made from cotton or synthetic materials, which help absorb sweat and keep the skin dry. However, in modern casting, fiberglass or polyester liners are often used, offering greater durability and comfort. Some liners are specifically designed to be waterproof, allowing patients to bathe or swim while wearing their casts. These waterproof liners are typically made from materials like polyurethane or special synthetic fibers that prevent water from seeping into the cast. While waterproof liners offer significant convenience, they may increase the application time and cost of the cast
Washable casts
There are some washable casts like FlexiOH which provide good ventilation and maintain good skin hygiene. With this cast, patients are able to bathe and go out in the rain. These types of casts have advantages that deliver patients a better treatment than conventional casts made of plaster of Paris or Fiberglass. They are the next generation of orthopedic immobilization photo-curing specialty-resin technology that enables a waterproof, washable, lightweight, strong and comfortable way of recovering from fractures.
Alternative methods of immobilization
Alternative immobilization techniques offer non-cast methods for stabilizing injuries, providing options that may be more comfortable, adjustable, or suitable for specific conditions. While traditional casts are commonly used for fractures and soft tissue injuries, alternatives are increasingly being utilized to address various patient needs and preferences.
One prominent alternative is the splint, which is often used for fractures that are not as severe or when a temporary immobilization method is required. Splints are typically made from materials like fiberglass, aluminum, or plastic and are easier to apply and adjust than casts. They can be used for injuries like sprains, minor fractures, or post-surgical stabilization. Unlike casts, splints are generally open on one or both sides, allowing for adjustments as swelling fluctuates during the healing process. They also provide more flexibility and can be removed for hygiene or rehabilitation purposes.
Orthopedic brace are another alternative, commonly used for joint injuries or soft tissue sprains and strains. Braces provide support and stabilize joints like the knee, ankle, or wrist. They are often used for conditions such as ligament sprains, tendinitis, or as post-operative support. Braces are usually made of fabric, neoprene, or metal components, allowing for greater mobility and easier removal compared to casts. They can be particularly useful for injuries that require gradual rehabilitation and controlled movement.
For certain types of fractures or injuries, Traction (orthopedics) is an effective immobilization method. Traction involves using a pulling force to align bones and reduce fractures, particularly in cases involving the spine, pelvis, or long bones. It can be achieved through a variety of mechanisms, including skin traction (using adhesive materials attached to the skin) or skeletal traction (which involves pins or wires placed directly into the bone). Traction helps maintain the correct alignment and promotes healing without the need for a cast, especially in more complex fractures.
In cases where the injury requires complete immobilization but not the rigidity of a cast, an orthopedic boot, also known as a CAM boot (controlled ankle motion) may be used, especially for foot or ankle injuries. These options are designed to protect the injured area while still allowing limited mobility. Orthopedic boots are often preferred in weight-bearing fractures, as they provide stability while allowing the patient to walk with crutches or other mobility aids.
Removal
Casts are typically removed by perforation using a cast saw, an oscillating saw designed to cut rigid material such as plaster or fiberglass while not harming soft tissue. Manually operated shears, patented in 1950 by Neil McKay, may be used on pediatric or other patients who may be affected by the noise of the saw.
History
The earliest methods of holding a reduced fracture involved using splints. These are rigid strips laid parallel to each other alongside the bone. The Ancient Egyptians used wooden splints made of bark wrapped in linen. They also used stiff bandages for support that were probably derived from embalming techniques. The use of plaster of Paris to cover walls is evident, but it seems it was never used for bandages. Ancient Hindus treated fractures with bamboo splints, and the writings of Hippocrates discuss management of fractures in some detail, recommending wooden splints plus exercise to prevent muscle atrophy during the immobilization. The ancient Greeks also used waxes and resins to create stiffened bandages and the Roman Celsus, writing in AD 30, describes how to use splints and bandages stiffened with starch. Arabian doctors used lime derived from sea shells and albumen from egg whites to stiffen bandages. The Italian School of Salerno in the twelfth century recommended bandages hardened with a flour and egg mixture as did medieval European bonesetters, who used casts made of egg white, flour, and animal fat. By the sixteenth century the famous French surgeon Ambroise Paré (1517–1590), who championed more humane treatments in medicine and promoted the use of artificial limbs, made casts of wax, cardboard, cloth, and parchment that hardened as they dried.
These methods all had merit, but the standard method for the healing of fractures was bed rest and restriction of activity. The search for a simpler, less-time-consuming, method led to the development of the first modern occlusive dressings, stiffened at first with starch and later with plaster-of-paris. The ambulatory treatment of fractures was the direct result of these innovations. The innovation of the modern cast can be traced to, among others, four military surgeons, Dominique Jean Larrey, Louis Seutin, Antonius Mathijsen, and Nikolai Ivanovich Pirogov.
Dominique Jean Larrey (1768–1842) was born in a small town in southern France. He first studied medicine with his uncle, a surgeon in Toulouse. After a short tour of duty as a naval surgeon, he returned to Paris, where he became caught up in the turmoil of the French Revolution, being present at the Storming of the Bastille. From then on, he made his career as a surgeon in France's revolutionary and Napoleonic armies, which he accompanied throughout Europe and the Middle East. As a result, Larrey accumulated a vast experience of military medicine and surgery.
One of his patients after the Battle of Borodino in 1812 was an infantry officer whose arm was amputated at the shoulder. The patient was evacuated immediately following the operation and passed from Russia, through Poland and Germany. When the dressing was removed on his arrival home in France, the wound had healed. Larrey concluded that the fact that the wound had been undisturbed had facilitated healing. After the war, Larrey began stiffening bandages using camphorated alcohol, lead acetate and egg whites beaten in water.
An improved method was introduced by Louis Seutin, (1793–1865) of Brussels. In 1815 Seutin had served in the allied armies in the war against Napoleon and was on the field of Waterloo. At the time of the development of his bandage he was chief surgeon in the Belgium army. Seutin's "bandage amidonnee" consisted of cardboard splints and bandages soaked in a solution of starch and applied wet. These dressings required 2 to 3 days to dry, depending on the temperature and humidity of the surroundings. The substitution of Dextrin for starch, advocated by Velpeau, the man widely regarded as the leading French surgeon at the beginning of the 19th century, reduced the drying time to 6 hours. Although this was a vast improvement, it was still a long time, especially in the harsh environment of the battlefield.
A good description of Seutin's technique was provided by Sampson Gamgee who learned it from Seutin in France during the winter of 1851–52 and went on to promote its use in Britain. The limb was initially wrapped in wool, especially over any bony prominences. Pasteboard was then cut into shape to provide a splint, and dampened so it could be molded to the limb. The limb was then wrapped in bandages before a starch coating was applied to the outer surface. Seutin's technique for the application of the starch apparatus formed the basis of the technique used with plaster of Paris dressings today. The use of this method led to the early mobilization of patients with fractures and a marked reduction in hospital time required.
Plaster casts
Although these bandages were an improvement over Larrey's method, they were far from ideal. They required a long time to apply and dry and there was often shrinkage and distortion. A great deal of interest had been aroused in Europe around 1800 by a British diplomat, consul William Eton, who described a method of treating fractures that he had observed in Turkey. He noted that gypsum plaster (plaster of Paris) was moulded around the patient's leg to cause immobilization. If the cast became loose due to atrophy or a reduction in swelling, then additional gypsum plaster was added to fill the space. Adapting the use of plaster of Paris for use in hospitals, however, took some time. In 1828, doctors in Berlin were treating leg fractures by aligning the bones in a long narrow box, which they filled with moist sand. Substitution of plaster of Paris for the sand was the next logical step. Such plaster casts did not succeed however as the patient was confined to bed due to the casts being heavy and cumbersome.
Plaster of Paris bandages were introduced in different forms by two army surgeons, one at a peacetime home station and another on active service at the front. Antonius Mathijsen (1805–1878) was born in Budel, the Netherlands, where his father was the village doctor. He was educated in Brussels, Maastricht and Utrecht obtaining the degree of doctor of medicine at Gissen in 1837. He spent his entire career as a medical officer in the Dutch Army. While he was stationed at Haarlem in 1851, he developed a method of applying plaster of Paris bandages. A brief note describing his method was published on January 30, 1852; it was followed shortly by more complete accounts. In these accounts Mathijsen emphasised that only simple materials were required and the bandage could be quickly applied without assistance. The bandages hardened rapidly, provided an exact fit and could be windowed or bivalved (cut to provide strain relief) easily. Mathijsen used coarsely woven materials, usually linen, into which dry plaster of Paris had been rubbed thoroughly. The bandages were then moistened with a wet sponge or brush as they were applied and rubbed by hand until they hardened.
Plaster of Paris dressings were first employed in the treatment of mass casualties in the 1850s during the Crimean War by Nikolai Ivanovich Pirogov (1810–1881). Pirogov was born in Moscow and received his early education there. After obtaining a medical degree at Dorpat (now Tartu, Estonia) he studied at Berlin and Göttingen before returning to Dorpat as a professor of Surgery. In 1840, he became the professor of surgery at the academy of military medicine in St. Petersburg. Pirogov introduced the use of ether anaesthesia to Russia and made important contributions to the study of cross-sectional human anatomy. With the help of his patron, the grand duchess Helene Pavlovna, he introduced female nurses into the military hospitals at the same time that Florence Nightingale was beginning a similar program in British military hospitals.
Seutin had travelled through Russia demonstrating his 'starched bandage', and his technique had been adopted by both the Russian army and navy by 1837. Pirogov had observed the use of plaster of Paris bandages in the studio of a sculptor who used strips of linen soaked in liquid plaster of Paris for making models (this technique, called "modroc," is still popular). Pirogov went on to develop his own methods, although he was aware of Mathijsen's work. Pirogov's method involved soaking coarse cloth in a plaster of Paris mixture immediately before application to the limbs, which were protected either by stockings or cotton pads. Large dressings were reinforced with pieces of wood.
As time passed and the method moved more into the mainstream some disagreement arose as to the problems associated with cutting off air to skin contact, and also some improvements were made. Eventually Pirogov's method gave way to Mathijsen's. Among the improvements suggested as early as 1860 was that of making the dressing resistant to water by painting the dried plaster of Paris with a mixture of shellac dissolved in alcohol. The first commercial bandages were not produced until 1931 in Germany, and were called Cellona. Before that the bandages were made by hand at the hospitals.
As a plaster cast is applied, it expands by approximately 0.5%. The less water used, the more linear expansion occurs. Potassium sulfate can be used as an accelerator and sodium borate as a retarder to control setting time.
| Technology | Devices | null |
2784391 | https://en.wikipedia.org/wiki/Cooksonia | Cooksonia | Cooksonia is an extinct group of primitive land plants, treated as a genus, although probably not monophyletic. The earliest Cooksonia date from the middle of the Silurian (the Wenlock epoch); the group continued to be an important component of the flora until the end of the Early Devonian, a total time span of . While Cooksonia fossils are distributed globally, most type specimens come from Britain, where they were first discovered in 1937. Cooksonia includes the oldest known plant to have a stem with vascular tissue and is thus a transitional form between the primitive non-vascular bryophytes and the vascular plants.
Description
Only the sporophyte phase of Cooksonia is currently known (i.e. the phase which produces spores rather than gametes). Individuals were small, a few centimetres tall, and had a simple structure. They lacked leaves, flowers and roots—although it has been speculated that they grew from a rhizome that has not been preserved. They had a simple stalk that branched dichotomously a few times. Each branch ended in a sporangium or spore-bearing capsule. In his original description of the genus, Lang described the sporangia as flattened, "with terminal sporangia that are short and wide", and in the species Cooksonia pertoni "considerably wider than high". A 2010 review of the genus by Gonez and Gerrienne produced a tighter definition, which requires the sporangia to be more-or-less trumpet-shaped (as in the illustration), with a 'lid' or operculum which disintegrates to release the spores.
Specimens of one species of Cooksonia have a dark stripe in the centre of their stalks, which has been interpreted as the earliest remains of water-carrying tissue. Other Cooksonia species lacked such conducting tissue.
Cooksonia specimens occur in a range of sizes, and vary in stem width from about 0.3 mm to 3 mm. Specimens of different sizes were probably different species, not fragments of larger organisms: fossils occur in consistent size groupings, and sporangia and spore details are different in organisms of different sizes. The organisms probably exhibited determinate growth (i.e. stems did not grow further after producing sporangia).
Some Cooksonia species bore stomata, which had a role in gas exchange; this was probably to assist in transpiration-driven transport of dissolved materials in the xylem, rather than primarily in photosynthesis, as suggested by their concentration at the tips of the axes. These clusterings of stomata are typically associated with a bulging in the axis at the neck of the sporangium, which may have contained photosynthetic tissue, reminiscent of some mosses.
As the genus is circumscribed by Gonez and Gerrienne, there are six possible species. C. pertoni, C. paranensis and C. banksii are all relatively similar with flat-topped, trumpet-shaped sporangia; stems are somewhat narrower in C. paranensis than in C. pertoni. Only one specimen of C. bohemica is known. It has stouter, more branched stems; the original shape of the sporangia is unclear because of poor preservation. C. hemisphaerica, described from the same locality as C. pertoni, differs in having sporangia of which the tops, at least as preserved, are hemispherical rather than flat. C. cambrensis also has spherical sporangia, but without the gradual widening at the base characteristic of the other species. Preservation of the sporangia is again poor. C. barrandei was described in 2018.
Physiology
While reconstructions traditionally depict Cooksonia as a green and red, photosynthesising, self-sufficient stem, it is likely that at least some fossils are of a sporophyte generation that was dependent on a gametophyte for its nutrition – a relationship that occurs in modern mosses and liverworts. However, no fossil evidence of a gametophyte of Cooksonia has been discovered to date.
The widths of Cooksonia fossils span an order of magnitude. Study of smaller Cooksonia fossils showed that once the tissue required to support the axes, protect them from desiccation, and transport water had been accounted for, no room remained for photosynthetic tissue, and the sporophyte may therefore have been dependent on the gametophyte. Further, the axis thickness is what would be expected if its sole role was to support a sporangium. It appears that, originally at least, the role of the axes in smaller species was solely to ensure continued spore dispersal, even if the axis desiccated. The potentially self-sufficient larger axes may represent the evolution of an independent sporophyte generation.
In 2018, the sporophyte of a new species, Cooksonia barrandei, was described, from about 432 million years ago. It is the oldest-known megafossil of land plants, . It was sufficiently robust to pass Boyce's test for possible self-sufficiency. Together with evidence that, unlike modern mosses and liverworts, hornwort sporophytes do have a degree of nutritional independence through photosynthesis, C. barrandei suggests that independent gametophyte and sporophyte generations could have been ancestral in land plants, rather than evolving later.
Taxonomy
The first Cooksonia species were described by William Henry Lang in 1937 and named in honor of Isabel Cookson, with whom he had collaborated and who collected specimens of Cooksonia pertoni in Perton Quarry, Wales, in 1934. There were originally two species, Cooksonia pertoni and C. hemisphaerica. The genus was defined as having narrow leafless stems (axes), which branched dichotomously, with terminal sporangia that were "short and wide". There was a central vascular cylinder consisting of annular tracheids (water-conducting cells with thickened walls). Six other species were later added to the genus: C. crassiparietilis, C. caledonica, C. cambrensis, C. bohemica, C. paranensis and C. banksii. A review in 2010 concluded that the delineation of the genus was inaccurate and that some species needed to be removed; in particular those in which sporangia were not more-or-less trumpet-shaped. As amended by Gonez and Gerrienne, Cooksonia has the following species:
C. pertoni Lang 1937 (the type species designated by Gonez & Gerrienne)
C. paranensis Gerrienne et al. 2001
Seven further species are considered doubtful because of the poor preservation of the specimens, but are left in the genus:
C. acuminata Mussa et al. 2002
C. barrandei Libertín et al. 2018
C. cambrensis Edwards 1979
C. degrezensis Senkevich
C. downtonensis Heard 1939
C. rusanovii Ananiev 1960
C. zhanyiensis Li & Cai 1978
Four species are excluded from the genus by Gonez and Gerrienne. Species that have been transferred or removed are:
C. banksii Habgood et al. 2002 now Concavatheca banksii (Habgood, Edwards & Axe 2002) Morris et al. 2012b
C. bohemica Schweitzer 1980 now Aberlemnia bohemica (Schweitzer 1980) Sakala, Pšenička & Kraft 2018
C. caledonica Edwards 1970 now Aberlemnia caledonica (Edwards 1970) Gonez & Gerrienne 2010
C. crassiparietilis Yurina 1964
C. hemisphaerica Lang 1937
C. caledonica and the less well-preserved C. crassiparietilis have sporangia which are composed of two 'valves', splitting to release their spores along a line opposite to where they are attached to the stem (i.e. distally).
Phylogeny
For some years, it was suspected that Cooksonia and its species were poorly characterized. Thus four different kinds of spore, probably representing four different species, were found in sporangia originally identified as C. pertoni.
A 2010 study of the genus produced the consensus cladogram shown below (some branches have been collapsed to reduce the size of the diagram). This was based on data from an earlier study (by Kenrick and Crane), supplemented by further information on Cooksonia species resulting from the authors' own research.
A more recent phylogeny by Hao and Xue from 2013:
This confirms that the genus Cooksonia sensu Lang (1937) is polyphyletic. A core group of five species are placed together, unresolved between the euphyllophytes and the lycophytes. The poorly preserved C. hemisphaerica is placed as the most basal tracheophyte. Two other species, C. crassiparietilis and C. caledonica, are placed in the stem group of the lycophytes. These two species have been removed from Cooksonia sensu Gonez & Gerrienne (C. caledonica has since been placed in a new genus Aberlemnia). Both have sporangia which, although borne terminally rather than laterally, have a mechanism for releasing spores similar to those of the zosterophylls.
A second cladistic analysis was carried out using only the three best preserved and thus best known species, C. pertoni, C. paranensis, and C. caledonica. The position of C. caledonica was confirmed, but C. pertoni and C. paranensis now formed a single clade more clearly related to the lycophytes than the euphyllophytes.
Cooksonioids
Cooksonia and similar genera have been placed in a group called "cooksonioids". Originally the term was used for a group of plants fitting the general description of Cooksonia (i.e. simple plants with naked axes showing dichotomous branching and terminal sporangia), but with uncertain evidence of vascular tissue. Boyce restricted the group to forms with axes usually less than 1 mm in diameter, and hence possibly not capable of independent growth. In addition to Cooksonia, he included genera such as Salopella, Tarrantia and Tortilicaulis. Hue and Xao regarded cooksonioids as a group within the rhyniophytes with radially symmetrical sporangia of roughly the same height and width, and included Cooksonia pertoni, C. paranensis and C. hemisphaerica, but not C. crassiparietilis and Aberlemnia caledonica, as they had bilaterally symmetrical sporangia.
| Biology and health sciences | Pteridophytes | Plants |
2784496 | https://en.wikipedia.org/wiki/Polyyne | Polyyne | A polyyne is any organic compound with alternating single and triple bonds; that is, a series of consecutive alkynes, with n greater than 1. These compounds are also called polyacetylenes, especially in the natural products and chemical ecology literature, even though this nomenclature more properly refers to acetylene polymers composed of alternating single and double bonds with n greater than 1. They are also sometimes referred to as oligoynes, or carbinoids after "carbyne" , the hypothetical allotrope of carbon that would be the ultimate member of the series. The synthesis of this substance has been claimed several times since the 1960s, but those reports have been disputed. Indeed, the substances identified as short chains of "carbyne" in many early organic synthesis attempts would be called polyynes today.
The simplest polyyne is diacetylene or butadiyne, . Along with cumulenes, polyynes are distinguished from other organic chains by their rigidity and high conductivity, both of which make them promising as wires in molecular nanotechnology. Polyynes have been detected in interstellar molecular clouds where hydrogen is scarce.
Synthesis
The first reported synthesis of a polyyne was performed in 1869 by , who observed that copper phenylacetylide () undergoes oxidative dimerization in the presence of air to produce diphenylbutadiyne ().
Interest in these compounds has stimulated research into their preparation by organic synthesis by several general routes. As a main synthetic tool usually acetylene homocoupling reactions like the Glaser coupling or its associated Elinton and Hay protocols are used. Moreover, many of such procedures involve a Cadiot–Chodkiewicz coupling or similar reactions to unite two separate alkyne building-blocks or by alkylation of a pre-formed polyyne unit. In addition to that, Fritsch–Buttenberg–Wiechell rearrangement was used as crucial step during the synthesis of the longest known polyyne (). An elimination of chlorovinylsilanes was used as a final step in the synthesis of the longest known phenyl end-capped polyynes.
Organic and organosilicon polyynes
Using various techniques, polyynes with n up to 4 or 5 were synthesized during the 1950s. Around 1971, T. R. Johnson and D. R. M. Walton developed the use of end-caps of the form –, where R was usually an ethyl group, to protect the polyyne chain during the chain-doubling reaction using Hay's catalyst (a copper(I)–TMEDA complex). With that technique they were able to obtain polyynes like with n up to 8 in pure state, and with n up to 16 in solution.
Later Tykwinski and co-workers were able to obtain polyynes with chain length up to C20.
A polyyne compound with 10 acetylenic units (20 atoms), with the ends capped by Fréchet-type aromatic polyether dendrimers, was isolated and characterized in 2002. Moreover, the synthesis of dicyanopolyynes with up to 8 acetylenic units was reported. The longest phenyl end-capped polyynes were reported by Cox and co-workers in 2007. As of 2010, the polyyne with the longest chain yet isolated had 22 acetylenic units (44 carbon atoms), end-capped with tris(3,5-di-t-butylphenyl)methyl groups.
Alkynes with the formula and n from 2 to 6 can be detected in the decomposition products of partially oxidized copper(I) acetylide ( (an acetylene derivative known since 1856 or earlier) by hydrochloric acid. A "carbonaceous" residue left by the decomposition also has the spectral signature of chains.
Organometallics
Organometallic polyynes capped with metal complexes are well characterized. As of the mid-2010s, the most intense research has concerned rhenium (, n = 3–10), ruthenium (, n = 4–10), iron (), platinum (, n = 8–14), palladium (, n = 3–5, Ar = aryl), and cobalt (, n = 7–13) complexes.
Stability
Long polyyne chains are said to be inherently unstable in bulk because they can cross-link with each other exothermically. Explosions are a real hazard in this area of research. They can be fairly stable, even against moisture and oxygen, if the end hydrogen atoms are replaced with a suitably inert end-group, such as tert-butyl or trifluoromethyl. Bulky end-groups, that can keep the chains apart, work especially well at stabilizing polyynes. In 1995 the preparation of carbyne chains with over 300 carbon atoms was reported using this technique. However the report has been contested by a claim that the detected molecules were fullerene-like structures rather than long polyynes.
Polyyne chains have also been stabilised to heating by co-deposition with silver nanoparticles, and by complexation with a mercury-containing tridentate Lewis acid to form layered adducts. Long polyyne chains encapsulated in double-walled carbon nanotubes or in the form of rotaxanes have also been shown to be stable. Despite rather low stability of longer polyynes there are some examples of their use as synthetic precursors in organic and organometallic synthesis.
Structure
Synthetic polyynes of the form , with n about 8 or more, often have a smoothly curved or helical backbone in the crystalline solid state, presumably due to crystal packing effects. For example, when the cap R is triisopropylsilyl and n is 8, X-ray crystallography of the substance (a crystalline orange/yellow solid) shows the backbone bent by about 25–30 degrees in a broad arch, so that each C−C≡C angle deviates by 3.1 degrees from a straight line. This geometry affords a denser packing, with the bulky cap of an adjacent molecule nested into the concave side of the backbone. As a result, the distance between backbones of neighboring molecules is reduced to about 0.35 to 0.5 nm, near the range at which one expects spontaneous cross-linking. The compound is stable indefinitely at low temperature, but decomposes before melting. In contrast, the homologous molecules with n = 4 or n = 5 have nearly straight backbones that stay at least 0.5 to 0.7 nm apart, and melt without decomposing.
Natural occurrence
Biological origins
A wide range of organisms synthesize polyynes. These chemicals have various biological activities, including as flavorings and pigments, chemical repellents and toxins, and potential application to biomedical research and pharmaceuticals. In plants, polyynes are found mainly in Asterids clade, especially in the sunflower, carrot, ginseng and bellflower families. However, they can also be found in some members of the tomato, olax, and sandalwood families. The earliest polyyne to be isolated was dehydromatricaria ester (DME) in 1826; however, it was not fully characterized until later.
The simple fatty acid 8,10-octadecadiynoic acid is isolated from the root bark of the legume Paramacrolobium coeruleum of the family Caesalpiniaceae and has been investigated as a photopolymerizable unit in synthetic phospholipids.
Thiarubrine B is the most prevalent among several related light-sensitive pigments that have been isolated from the Giant Ragweed (Ambrosia trifida), a plant used in herbal medicine. The thiarubrines have antibiotic, antiviral, and nematocidal activity, and activity against HIV-1 that is mediated by exposure to light.
Polyynes such as falcarindiol can be found in Apiaceae vegetables like carrot, celery, fennel, parsley and parsnip where they show cytotoxic activities. Aliphatic -polyynes of the falcarinol type were described to act as metabolic modulators and are studied as potential health-promoting nutraceuticals. Falcarindiol is the main compound responsible for bitterness in carrots, and is the most active among several polyynes with potential anticancer activity found in Devil's club (Oplopanax horridus). Other polyynes from plants include oenanthotoxin and cicutoxin, which are poisons found in water dropwort (Oenanthe spp.) and water hemlock (Cicuta spp.).
Ichthyothere is a genus of plants whose active constituent is a polyyne called ichthyothereol. This compound is highly toxic to fish and mammals. Ichthyothere terminalis leaves have traditionally been used to make poisoned bait by indigenous peoples of the lower Amazon basin.
Dihydromatricaria acid is a polyyne produced and secreted by soldier beetles as a chemical defense.
In space
The octatetraynyl radicals and hexatriynyl radicals together with their ions are detected in space where hydrogen is rare. Moreover, there have been claims that polyynes have been found in astronomical impact sites on Earth as part of the mineral chaoite, but this interpretation has been contested. See Astrochemistry.
| Physical sciences | Aliphatic hydrocarbons | Chemistry |
1364502 | https://en.wikipedia.org/wiki/Belt%20%28mechanical%29 | Belt (mechanical) | A belt is a loop of flexible material used to link two or more rotating shafts mechanically, most often parallel. Belts may be used as a source of motion, to transmit power efficiently or to track relative movement. Belts are looped over pulleys and may have a twist between the pulleys, and the shafts need not be parallel.
In a two pulley system, the belt can either drive the pulleys normally in one direction (the same if on parallel shafts), or the belt may be crossed, so that the direction of the driven shaft is reversed (the opposite direction to the driver if on parallel shafts). The belt drive can also be used to change the speed of rotation, either up or down, by using different sized pulleys.
As a source of motion, a conveyor belt is one application where the belt is adapted to carry a load continuously between two points.
History
The mechanical belt drive, using a pulley machine, was first mentioned in the text of the Dictionary of Local Expressions by the Han Dynasty philosopher, poet, and politician Yang Xiong (53–18 BC) in 15 BC, used for a quilling machine that wound silk fibres onto bobbins for weavers' shuttles. The belt drive is an essential component of the invention of the spinning wheel. The belt drive was not only used in textile technologies, it was also applied to hydraulic-powered bellows dated from the 1st century AD.
Power transmission
Belts are the cheapest utility for power transmission between shafts that may not be axially aligned. Power transmission is achieved by purposely designed belts and pulleys. The variety of power transmission needs that can be met by a belt-drive transmission system are numerous, and this has led to many variations on the theme. Belt drives run smoothly and with little noise, and provide shock absorption for motors, loads, and bearings when the force and power needed changes. A drawback to belt drives is that they transmit less power than gears or chain drives. However, improvements in belt engineering allow use of belts in systems that formerly only allowed chain drives or gears.
Power transmitted between a belt and a pulley is expressed as the product of difference of tension and belt velocity:
where and are tensions in the tight side and slack side of the belt respectively. They are related as
where is the coefficient of friction, and is the angle (in radians) subtended by contact surface at the centre of the pulley.
Power transmission loss form
Pros and cons
Belt drives are simple, inexpensive, and do not require axially aligned shafts. They help protect machinery from overload and jam, and damp and isolate noise and vibration. Load fluctuations are shock-absorbed (cushioned). They need no lubrication and minimal maintenance. They have high efficiency (90–98%, usually 95%), high tolerance for misalignment, and are of relatively low cost if the shafts are far apart. Clutch action can be achieved by shifting the belt to a free turning pulley or by releasing belt tension. Different speeds can be obtained by stepped or tapered pulleys.
The angular-velocity ratio may not be exactly constant or equal to that of the pulley diameters, due to slip and stretch. However, this problem can be largely solved by the use of toothed belts. Working temperatures range from . Adjustment of centre distance or addition of an idler pulley is crucial to compensate for wear and stretch.
Flat belts
Flat belts were widely used in the 19th and early 20th centuries in line shafting to transmit power in factories. They were also used in countless farming, mining, and logging applications, such as bucksaws, sawmills, threshers, silo blowers, conveyors for filling corn cribs or haylofts, balers, water pumps (for wells, mines, or swampy farm fields), and electrical generators. Flat belts are still used today, although not nearly as much as in the line-shaft era. The flat belt is a simple system of power transmission that was well suited for its day. It can deliver high power at high speeds (373 kW at 51 m/s; 115 mph), in cases of wide belts and large pulleys. Wide-belt-large-pulley drives are bulky, consuming much space while requiring high tension, leading to high loads, and are poorly suited to close-centers applications. V-belts have mainly replaced flat belts for short-distance power transmission; and longer-distance power transmission is typically no longer done with belts at all. For example, factory machines now tend to have individual electric motors.
Because flat belts tend to climb towards the higher side of the pulley, pulleys were made with a slightly convex or "crowned" surface (rather than flat) to allow the belt to self-center as it runs. Flat belts also tend to slip on the pulley face when heavy loads are applied, and many proprietary belt dressings were available that could be applied to the belts to increase friction, and so power transmission.
Flat belts were traditionally made of leather or fabric. Early flour mills in Ukraine had leather belt drives. After World War I, there was such a shortage of shoe leather that people cut up the belt drives to make shoes. Selling shoes was more profitable than selling flour for a time. Flour milling soon came to a standstill and bread prices rose, contributing to famine conditions. Leather drive belts were put to another use during the Rhodesian Bush War (1964–1979): To protect riders of cars and busses from land mines, layers of leather belt drives were placed on the floors of vehicles in danger zones. Today most belt drives are made of rubber or synthetic polymers. Grip of leather belts is often better if they are assembled with the hair side (outer side) of the leather against the pulley, although some belts are instead given a half-twist before joining the ends (forming a Möbius strip), so that wear can be evenly distributed on both sides of the belt. Belts ends are joined by lacing the ends together with leather thonging (the oldest of the methods), steel comb fasteners and/or lacing, or by gluing or welding (in the case of polyurethane or polyester). Flat belts were traditionally jointed, and still usually are, but they can also be made with endless construction.
Rope drives
In the mid 19th century, British millwrights discovered that multi-grooved pulleys connected by ropes outperformed flat pulleys connected by leather belts. Wire ropes were occasionally used, but cotton, hemp, manila hemp and flax rope saw the widest use. Typically, the rope connecting two pulleys with multiple V-grooves was spliced into a single loop that traveled along a helical path before being returned to its starting position by an idler pulley that also served to maintain the tension on the rope. Sometimes, a single rope was used to transfer power from one multiple-groove drive pulley to several single- or multiple-groove driven pulleys in this way.
In general, as with flat belts, rope drives were used for connections from stationary engines to the jack shafts and line shafts of mills, and sometimes from line shafts to driven machinery. Unlike leather belts, however, rope drives were sometimes used to transmit power over relatively long distances. Over long distances, intermediate sheaves were used to support the "flying rope", and in the late 19th century, this was considered quite efficient.
Round belts
Round belts are a circular cross section belt designed to run in a pulley with a 60 degree V-groove. Round grooves are only suitable for idler pulleys that guide the belt, or when (soft) O-ring type belts are used. The V-groove transmits torque through a wedging action, thus increasing friction. Nevertheless, round belts are for use in relatively low torque situations only and may be purchased in various lengths or cut to length and joined, either by a staple, a metallic connector (in the case of hollow plastic), gluing or welding (in the case of polyurethane). Early sewing machines utilized a leather belt, joined either by a metal staple or glued, to great effect.
Spring belts
Spring belts are similar to rope or round belts but consist of a long steel helical spring. They are commonly found on toy or small model engines, typically steam engines driving other toys or models or providing a transmission between the crankshaft and other parts of a vehicle. The main advantage over rubber or other elastic belts is that they last much longer under poorly controlled operating conditions. The distance between the pulleys is also less critical. Their main disadvantage is that slippage is more likely due to the lower coefficient of friction. The ends of a spring belt can be joined either by bending the last turn of the helix at each end by 90 degrees to form hooks, or by reducing the diameter of the last few turns at one end so that it "screws" into the other end.
V belts (also style V-belts, vee belts, or, less commonly, wedge rope) solved the slippage and alignment problem. It is now the basic belt for power transmission. They provide the best combination of traction, speed of movement, load of the bearings, and long service life. They are generally endless, and their general cross-section shape is roughly trapezoidal (hence the name "V"). The "V" shape of the belt tracks in a mating groove in the pulley (or sheave), with the result that the belt cannot slip off. The belt also tends to wedge into the groove as the load increases—the greater the load, the greater the wedging action—improving torque transmission and making the V-belt an effective solution, needing less width and tension than flat belts. V-belts trump flat belts with their small center distances and high reduction ratios. The preferred center distance is larger than the largest pulley diameter, but less than three times the sum of both pulleys. Optimal speed range is . V-belts need larger pulleys for their thicker cross-section than flat belts.
For high-power requirements, two or more V-belts can be joined side-by-side in an arrangement called a multi-V, running on matching multi-groove sheaves. This is known as a multiple-V-belt drive (or sometimes a "classical V-belt drive").
V-belts may be homogeneously rubber or polymer throughout, or there may be fibers embedded in the rubber or polymer for strength and reinforcement. The fibers may be of textile materials such as cotton, polyamide (such as nylon) or polyester or, for greatest strength, of steel or aramid (such as Technora, Twaron or Kevlar).
When an endless belt does not fit the need, jointed and link V-belts may be employed. Most models offer the same power and speed ratings as equivalently-sized endless belts and do not require special pulleys to operate. A link v-belt is a number of polyurethane/polyester composite links held together, either by themselves, such as Fenner Drives' PowerTwist, or Nu-T-Link (with metal studs). These provide easy installation and superior environmental resistance compared to rubber belts and are length-adjustable by disassembling and removing links when needed.
History of V-belts
Trade journal coverage of V-belts in automobiles from 1916 mentioned leather as the belt material, and mentioned that the V angle was not yet well standardized. The endless rubber V-belt was developed in 1917 by Charles C. Gates of the Gates Rubber Company. Multiple-V-belt drive was first arranged a few years later by Walter Geist of the Allis-Chalmers corporation, who was inspired to replace the single rope of multi-groove-sheave rope drives with multiple V-belts running parallel. Geist filed for a patent in 1925, and Allis-Chalmers began marketing the drive under the "Texrope" brand; the patent was granted in 1928 (). The "Texrope" brand still exists, although it has changed ownership and no longer refers to multiple-V-belt drive alone.
Multi-groove belts
A multi-groove, V-ribbed, or polygroove belt is made up of usually between 3 and 24 V-shaped sections alongside each other. This gives a thinner belt for the same drive surface, thus it is more flexible, although often wider. The added flexibility offers an improved efficiency, as less energy is wasted in the internal friction of continually bending the belt. In practice this gain of efficiency causes a reduced heating effect on the belt, and a cooler-running belt lasts longer in service. Belts are commercially available in several sizes, with usually a 'P' (sometimes omitted) and a single letter identifying the pitch between grooves. The 'PK' section with a pitch of 3.56 mm is commonly used for automotive applications.
A further advantage of the polygroove belt that makes them popular is that they can run over pulleys on the ungrooved back of the belt. Though this is sometimes done with V-belts with a single idler pulley for tensioning, a polygroove belt may be wrapped around a pulley on its back tightly enough to change its direction, or even to provide a light driving force.
Any V-belt's ability to drive pulleys depends on wrapping the belt around a sufficient angle of the pulley to provide grip. Where a single-V-belt is limited to a simple convex shape, it can adequately wrap at most three or possibly four pulleys, so can drive at most three accessories. Where more must be driven, such as for modern cars with power steering and air conditioning, multiple belts are required. As the polygroove belt can be bent into concave paths by external idlers, it can wrap any number of driven pulleys, limited only by the power capacity of the belt.
This ability to bend the belt at the designer's whim allows it to take a complex or "serpentine" path. This can assist the design of a compact engine layout, where the accessories are mounted more closely to the engine block and without the need to provide movable tensioning adjustments. The entire belt may be tensioned by a single idler pulley.
The nomenclature used for belt sizes varies by region and trade. An automotive belt with the number "740K6" or "6K740" indicates a belt in length, 6 ribs wide, with a rib pitch of (a standard thickness for a K series automotive belt would be 4.5mm). A metric equivalent would be usually indicated by "6PK1880" whereby 6 refers to the number of ribs, PK refers to the metric PK thickness and pitch standard, and 1880 is the length of the belt in millimeters.
Ribbed belt
A ribbed belt is a power transmission belt featuring lengthwise grooves. It operates from contact between the ribs of the belt and the grooves in the pulley. Its single-piece structure is reported to offer an even distribution of tension across the width of the pulley where the belt is in contact, a power range up to 600 kW, a high speed ratio, serpentine drives (possibility to drive off the back of the belt), long life, stability and homogeneity of the drive tension, and reduced vibration. The ribbed belt may be fitted on various applications: compressors, fitness bikes, agricultural machinery, food mixers, washing machines, lawn mowers, etc.
Film belts
Though often grouped with flat belts, they are actually a different kind. They consist of a very thin belt (0.5–15 millimeters or 100–4000 micrometres) strip of plastic and occasionally rubber. They are generally intended for low-power (less than 10 watts), high-speed uses, allowing high efficiency (up to 98%) and long life. These are seen in business machines, printers, tape recorders, and other light-duty operations.
Timing belts
Timing belts (also known as toothed, notch, cog, or synchronous belts) are a positive transfer belt and can track relative movement. These belts have teeth that fit into a matching toothed pulley. When correctly tensioned, they have no slippage, run at constant speed, and are often used to transfer direct motion for indexing or timing purposes (hence their name). They are often used instead of chains or gears, so there is less noise and a lubrication bath is not necessary. Camshafts of automobiles, miniature timing systems, and stepper motors often utilize these belts. Timing belts need the least tension of all belts and are among the most efficient. They can bear up to at speeds of .
Timing belts with a helical offset tooth design are available. The helical offset tooth design forms a chevron pattern and causes the teeth to engage progressively. The chevron pattern design is self-aligning and does not make the noise that some timing belts make at certain speeds, and is more efficient at transferring power (up to 98%).
The advantages of timing belts include clean operation, energy efficiency, low maintenance, low noise, non slip performance, versatile load and speed capabilities.
Disadvantages include a relatively high purchase cost, the need for specially fabricated toothed pulleys, less protection from overloading, jamming, and vibration due to their continuous tension cords, the lack of clutch action (only possible with friction-drive belts), and the fixed lengths, which do not allow length adjustment (unlike link V-belts or chains).
Specialty belts
Belts normally transmit power on the tension side of the loop. However, designs for continuously variable transmissions exist that use belts that are a series of solid metal blocks, linked together as in a chain, transmitting power on the compression side of the loop.
Rolling roads
Belts used for rolling roads for wind tunnels can be capable of .
Standards for use
The open belt drive has parallel shafts rotating in the same direction, whereas the cross-belt drive also bears parallel shafts but rotate in opposite direction. The former is far more common, and the latter not appropriate for timing and standard V-belts unless there is a twist between each pulley so that the pulleys only contact the same belt surface. Nonparallel shafts can be connected if the belt's center line is aligned with the center plane of the pulley. Industrial belts are usually reinforced rubber but sometimes leather types. Non-leather, non-reinforced belts can only be used in light applications.
The pitch line is the line between the inner and outer surfaces that is neither subject to tension (like the outer surface) nor compression (like the inner). It is midway through the surfaces in film and flat belts and dependent on cross-sectional shape and size in timing and V-belts. Standard reference pitch diameter can be estimated by taking average of gear teeth tips diameter and gear teeth base diameter. The angular speed is inversely proportional to size, so the larger the one wheel, the less angular velocity, and vice versa. Actual pulley speeds tend to be 0.5–1% less than generally calculated because of belt slip and stretch. In timing belts, the inverse ratio teeth of the belt contributes to the exact measurement.
The speed of the belt is:
International use standards
Standards include:
ISO 9563: This standard specifies requirements and test methods for endless power transmission V-belts and V-ribbed belts.
ISO 4184: This standard specifies the dimensions of classical and narrow V-belts for general use.
ISO 9981: This standard deals with the dimensions of rubber synchronous belt drives.
ISO 9982: This standard covers the dimensions of polyurethane synchronous belt drives.
DIN 22101: This standard covers the design principles for belt conveyors used in bulk material handling, including safety requirements and testing methods.
ASME B29.1: This standard specifies the dimensions, tolerances, and quality requirements for roller chain drives, which include belts and sprockets.
ANSI/RMA IP-20 is a standard developed by the American National Standards Institute (ANSI) and the Rubber Manufacturers Association (RMA) that focuses on elastomeric belts used in industrial applications. This standard covers important aspects such as dimensions and tolerances, ensuring that the belts perform reliably and efficiently in various industrial settings.
SAE J1459 is a standard developed by the Society of Automotive Engineers (SAE) that focuses on automotive V-belts and V-ribbed belts. These belts are used in various automotive applications, such as power transmission between the engine and different accessories, including the alternator, power steering pump, air conditioning compressor, and water pump. The standard specifies test procedures, performance requirements, and dimensions to ensure the belts are reliable, durable, and suitable for automotive use.
ASTM D378 is a standard developed by the American Society for Testing and Materials (ASTM), which focuses on the testing of conveyor belts used in various industries for specific applications. Conveyor belts are essential for material handling and transportation in industries such as mining, construction, agriculture, and manufacturing. ASTM D378 covers the testing methods to evaluate conveyor belts for performance characteristics, such as fire resistance and oil resistance, ensuring that they meet safety and operational requirements.
Selection criteria
Belt drives are built under the following required conditions: speeds of and power transmitted between drive and driven unit; suitable distance between shafts; and appropriate operating conditions. The equation for power is
Factors of power adjustment include speed ratio; shaft distance (long or short); type of drive unit (electric motor, internal combustion engine); service environment (oily, wet, dusty); driven unit loads (jerky, shock, reversed); and pulley-belt arrangement (open, crossed, turned). These are found in engineering handbooks and manufacturer's literature. When corrected, the power is compared to rated powers of the standard belt cross-sections at particular belt speeds to find a number of arrays that perform best. Now the pulley diameters are chosen. It is generally either large diameters or large cross-section that are chosen, since, as stated earlier, larger belts transmit this same power at low belt speeds as smaller belts do at high speeds. To keep the driving part at its smallest, minimal-diameter pulleys are desired. Minimum pulley diameters are limited by the elongation of the belt's outer fibers as the belt wraps around the pulleys. Small pulleys increase this elongation, greatly reducing belt life. Minimal pulley diameters are often listed with each cross-section and speed, or listed separately by belt cross-section. After the cheapest diameters and belt section are chosen, the belt length is computed. If endless belts are used, the desired shaft spacing may need adjusting to accommodate standard-length belts. It is often more economical to use two or more juxtaposed V-belts, rather than one larger belt.
In large speed ratios or small central distances, the angle of contact between the belt and pulley may be less than 180°. If this is the case, the drive power must be further increased, according to manufacturer's tables, and the selection process repeated. This is because power capacities are based on the standard of a 180° contact angle. Smaller contact angles mean less area for the belt to obtain traction, and thus the belt carries less power.
Belt friction
Belt drives depend on friction to operate, but excessive friction wastes energy and rapidly wears the belt. Factors that affect belt friction include belt tension, contact angle, and the materials used to make the belt and pulleys.
Belt tension
Power transmission is a function of belt tension. However, also increasing with tension is stress (load) on the belt and bearings. The ideal belt is that of the lowest tension that does not slip in high loads. Belt tensions should also be adjusted to belt type, size, speed, and pulley diameters. Belt tension is determined by measuring the force to deflect the belt a given distance per inch (or mm) of pulley. Timing belts need only adequate tension to keep the belt in contact with the pulley.
Belt wear
Fatigue, more so than abrasion, is the culprit for most belt problems. This wear is caused by stress from rolling around the pulleys. High belt tension; excessive slippage; adverse environmental conditions; and belt overloads caused by shock, vibration, or belt slapping all contribute to belt fatigue.
Belt vibration
Vibration signatures are widely used for studying belt drive malfunctions. Some of the common malfunctions or faults include the effects of belt tension, speed, sheave eccentricity and misalignment conditions. The effect of sheave Eccentricity on vibration signatures of the belt drive is quite significant. Although, vibration magnitude is not necessarily increased by this it will create strong amplitude modulation. When the top section of a belt is in resonance, the vibrations of the machine is increased. However, an increase in the machine vibration is not significant when only the bottom section of the belt is in resonance. The vibration spectrum has the tendency to move to higher frequencies as the tension force of the belt is increased.
Belt dressing
Belt slippage can be addressed in several ways. Belt replacement is an obvious solution, and eventually the mandatory one (because no belt lasts forever). Often, though, before the replacement option is executed, retensioning (via pulley centerline adjustment) or dressing (with any of various coatings) may be successful to extend the belt's lifespan and postpone replacement. Belt dressings are typically liquids that are poured, brushed, dripped, or sprayed onto the belt surface and allowed to spread around; they are meant to recondition the belt's driving surfaces and increase friction between the belt and the pulleys. Some belt dressings are dark and sticky, resembling tar or syrup; some are thin and clear, resembling mineral spirits. Some are sold to the public in aerosol cans at auto parts stores; others are sold in drums only to industrial users.
Specifications
To fully specify a belt, the material, length, and cross-section size and shape are required. Timing belts, in addition, require that the size of the teeth be given. The length of the belt is the sum of the central length of the system on both sides, half the circumference of both pulleys, and the square of the sum (if crossed) or the difference (if open) of the radii. Thus, when dividing by the central distance, it can be visualized as the central distance times the height that gives the same squared value of the radius difference on, of course, both sides. When adding to the length of either side, the length of the belt increases, in a similar manner to the Pythagorean theorem. One important concept to remember is that as gets closer to there is less of a distance (and therefore less addition of length) as it approaches zero.
On the other hand, in a crossed belt drive the sum rather than the difference of radii is the basis for computation for length. So the wider the small drive increases, the belt length is higher.
V-belt profiles
Metric v-belt profiles (note pulley angles are reduced for small radius pulleys):
* Common pulley design is to have a higher angle of the first part of the opening, above the so-called "pitch line".
E.g. the pitch line for SPZ could be 8.5 mm from the bottom of the "V". In other words, 0–8.5 mm is 35° and 45° from 8.5 and above.
| Technology | Flexible components | null |
1364622 | https://en.wikipedia.org/wiki/Four-dimensional%20space | Four-dimensional space | Four-dimensional space (4D) is the mathematical extension of the concept of three-dimensional space (3D). Three-dimensional space is the simplest possible abstraction of the observation that one needs only three numbers, called dimensions, to describe the sizes or locations of objects in the everyday world. For example, the volume of a rectangular box is found by measuring and multiplying its length, width, and height (often labeled , , and ). This concept of ordinary space is called Euclidean space because it corresponds to Euclid's geometry, which was originally abstracted from the spatial experiences of everyday life.
The idea of adding a fourth dimension appears in Jean le Rond d'Alembert's "Dimensions", published in 1754, but the mathematics of more than three dimensions only emerged in the 19th century. The general concept of Euclidean space with any number of dimensions was fully developed by the Swiss mathematician Ludwig Schläfli before 1853. Schläfli's work received little attention during his lifetime and was published only posthumously, in 1901, but meanwhile the fourth Euclidean dimension was rediscovered by others. In 1880 Charles Howard Hinton popularized it in an essay, "What is the Fourth Dimension?", in which he explained the concept of a "four-dimensional cube" with a step-by-step generalization of the properties of lines, squares, and cubes. The simplest form of Hinton's method is to draw two ordinary 3D cubes in 2D space, one encompassing the other, separated by an "unseen" distance, and then draw lines between their equivalent vertices. This can be seen in the accompanying animation whenever it shows a smaller inner cube inside a larger outer cube. The eight lines connecting the vertices of the two cubes in this case represent a single direction in the "unseen" fourth dimension.
Higher-dimensional spaces (greater than three) have since become one of the foundations for formally expressing modern mathematics and physics. Large parts of these topics could not exist in their current forms without using such spaces. Einstein's theory of relativity is formulated in 4D space, although not in a Euclidean 4D space. Einstein's concept of spacetime has a Minkowski structure based on a non-Euclidean geometry with three spatial dimensions and one temporal dimension, rather than the four symmetric spatial dimensions of Schläfli's Euclidean 4D space.
Single locations in Euclidean 4D space can be given as vectors or 4-tuples, i.e., as ordered lists of numbers such as . It is only when such locations are linked together into more complicated shapes that the full richness and geometric complexity of higher-dimensional spaces emerge. A hint of that complexity can be seen in the accompanying 2D animation of one of the simplest possible regular 4D objects, the tesseract, which is analogous to the 3D cube.
History
Lagrange wrote in his (published 1788, based on work done around 1755) that mechanics can be viewed as operating in a four-dimensional space— three dimensions of space, and one of time. As early as 1827, Möbius realized that a fourth spatial dimension would allow a three-dimensional form to be rotated onto its mirror-image. The general concept of Euclidean space with any number of dimensions was fully developed by the Swiss mathematician Ludwig Schläfli in the mid-19th century, at a time when Cayley, Grassman and Möbius were the only other people who had ever conceived the possibility of geometry in more than three dimensions. By 1853 Schläfli had discovered all the regular polytopes that exist in higher dimensions, including the four-dimensional analogs of the Platonic solids.
An arithmetic of four spatial dimensions, called quaternions, was defined by William Rowan Hamilton in 1843. This associative algebra was the source of the science of vector analysis in three dimensions as recounted by Michael J. Crowe in A History of Vector Analysis. Soon after, tessarines and coquaternions were introduced as other four-dimensional algebras over R. In 1886, Victor Schlegel described his method of visualizing four-dimensional objects with Schlegel diagrams.
One of the first popular expositors of the fourth dimension was Charles Howard Hinton, starting in 1880 with his essay What is the Fourth Dimension?, published in the Dublin University magazine. He coined the terms tesseract, ana and kata in his book A New Era of Thought and introduced a method for visualizing the fourth dimension using cubes in the book Fourth Dimension. Hinton's ideas inspired a fantasy about a "Church of the Fourth Dimension" featured by Martin Gardner in his January 1962 "Mathematical Games column" in Scientific American.
Higher dimensional non-Euclidean spaces were put on a firm footing by Bernhard Riemann's 1854 thesis, , in which he considered a "point" to be any sequence of coordinates . In 1908, Hermann Minkowski presented a paper consolidating the role of time as the fourth dimension of spacetime, the basis for Einstein's theories of special and general relativity. But the geometry of spacetime, being non-Euclidean, is profoundly different from that explored by Schläfli and popularised by Hinton. The study of Minkowski space required Riemann's mathematics which is quite different from that of four-dimensional Euclidean space, and so developed along quite different lines. This separation was less clear in the popular imagination, with works of fiction and philosophy blurring the distinction, so in 1973 H. S. M. Coxeter felt compelled to write:
Vectors
Mathematically, a four-dimensional space is a space that needs four parameters to specify a point in it. For example, a general point might have position vector , equal to
This can be written in terms of the four standard basis vectors , given by
so the general vector is
Vectors add, subtract and scale as in three dimensions.
The dot product of Euclidean three-dimensional space generalizes to four dimensions as
It can be used to calculate the norm or length of a vector,
and calculate or define the angle between two non-zero vectors as
Minkowski spacetime is four-dimensional space with geometry defined by a non-degenerate pairing different from the dot product:
As an example, the distance squared between the points and is 3 in both the Euclidean and Minkowskian 4-spaces, while the distance squared between and is 4 in Euclidean space and 2 in Minkowski space; increasing decreases the metric distance. This leads to many of the well-known apparent "paradoxes" of relativity.
The cross product is not defined in four dimensions. Instead, the exterior product is used for some applications, and is defined as follows:
This is bivector valued, with bivectors in four dimensions forming a six-dimensional linear space with basis . They can be used to generate rotations in four dimensions.
Orthogonality and vocabulary
In the familiar three-dimensional space of daily life, there are three coordinate axes—usually labeled , , and —with each axis orthogonal (i.e. perpendicular) to the other two. The six cardinal directions in this space can be called up, down, east, west, north, and south. Positions along these axes can be called altitude, longitude, and latitude. Lengths measured along these axes can be called height, width, and depth.
Comparatively, four-dimensional space has an extra coordinate axis, orthogonal to the other three, which is usually labeled . To describe the two additional cardinal directions, Charles Howard Hinton coined the terms ana and kata, from the Greek words meaning "up toward" and "down from", respectively.
As mentioned above, Hermann Minkowski exploited the idea of four dimensions to discuss cosmology including the finite velocity of light. In appending a time dimension to three-dimensional space, he specified an alternative perpendicularity, hyperbolic orthogonality. This notion provides his four-dimensional space with a modified simultaneity appropriate to electromagnetic relations in his cosmos. Minkowski's world overcame problems associated with the traditional absolute space and time cosmology previously used in a universe of three space dimensions and one time dimension.
Geometry
The geometry of four-dimensional space is much more complex than that of three-dimensional space, due to the extra degree of freedom.
Just as in three dimensions there are polyhedra made of two dimensional polygons, in four dimensions there are polychora made of polyhedra. In three dimensions, there are 5 regular polyhedra known as the Platonic solids. In four dimensions, there are 6 convex regular 4-polytopes, the analogs of the Platonic solids. Relaxing the conditions for regularity generates a further 58 convex uniform 4-polytopes, analogous to the 13 semi-regular Archimedean solids in three dimensions. Relaxing the conditions for convexity generates a further 10 nonconvex regular 4-polytopes.
In three dimensions, a circle may be extruded to form a cylinder. In four dimensions, there are several different cylinder-like objects. A sphere may be extruded to obtain a spherical cylinder (a cylinder with spherical "caps", known as a spherinder), and a cylinder may be extruded to obtain a cylindrical prism (a cubinder). The Cartesian product of two circles may be taken to obtain a duocylinder. All three can "roll" in four-dimensional space, each with its properties.
In three dimensions, curves can form knots but surfaces cannot (unless they are self-intersecting). In four dimensions, however, knots made using curves can be trivially untied by displacing them in the fourth direction—but 2D surfaces can form non-trivial, non-self-intersecting knots in 4D space. Because these surfaces are two-dimensional, they can form much more complex knots than strings in 3D space can. The Klein bottle is an example of such a knotted surface. Another such surface is the real projective plane.
Hypersphere
The set of points in Euclidean 4-space having the same distance from a fixed point forms a hypersurface known as a 3-sphere. The hyper-volume of the enclosed space is:
This is part of the Friedmann–Lemaître–Robertson–Walker metric in General relativity where is substituted by function with meaning the cosmological age of the universe. Growing or shrinking with time means expanding or collapsing universe, depending on the mass density inside.
Four-dimensional perception in humans
Research using virtual reality finds that humans, despite living in a three-dimensional world, can, without special practice, make spatial judgments about line segments embedded in four-dimensional space, based on their length (one-dimensional) and the angle (two-dimensional) between them. The researchers noted that "the participants in our study had minimal practice in these tasks, and it remains an open question whether it is possible to obtain more sustainable, definitive, and richer 4D representations with increased perceptual experience in 4D virtual environments". In another study, the ability of humans to orient themselves in 2D, 3D, and 4D mazes has been tested. Each maze consisted of four path segments of random length and connected with orthogonal random bends, but without branches or loops (i.e. actually labyrinths). The graphical interface was based on John McIntosh's free 4D Maze game. The participating persons had to navigate through the path and finally estimate the linear direction back to the starting point. The researchers found that some of the participants were able to mentally integrate their path after some practice in 4D (the lower-dimensional cases were for comparison and for the participants to learn the method).
However, a 2020 review underlined how these studies are composed of a small subject sample and mainly of college students. It also pointed out other issues that future research has to resolve: elimination of artifacts (these could be caused, for example, by strategies to resolve the required task that don't use 4D representation/4D reasoning and feedback given by researchers to speed up the adaptation process) and analysis on inter-subject variability (if 4D perception is possible, its acquisition could be limited to a subset of humans, to a specific critical period, or to people's attention or motivation). Furthermore, it is undetermined if there is a more appropriate way to project the 4-dimension (because there are no restrictions on how the 4-dimension can be projected). Researchers also hypothesized that human acquisition of 4D perception could result in the activation of brain visual areas and entorhinal cortex. If so they suggest that it could be used as a strong indicator of 4D space perception acquisition. Authors also suggested using a variety of different neural network architectures (with different a priori assumptions) to understand which ones are or are not able to learn.
Dimensional analogy
To understand the nature of four-dimensional space, a device called dimensional analogy is commonly employed. Dimensional analogy is the study of how () dimensions relate to dimensions, and then inferring how dimensions would relate to () dimensions.
The dimensional analogy was used by Edwin Abbott Abbott in the book Flatland, which narrates a story about a square that lives in a two-dimensional world, like the surface of a piece of paper. From the perspective of this square, a three-dimensional being has seemingly god-like powers, such as ability to remove objects from a safe without breaking it open (by moving them across the third dimension), to see everything that from the two-dimensional perspective is enclosed behind walls, and to remain completely invisible by standing a few inches away in the third dimension.
By applying dimensional analogy, one can infer that a four-dimensional being would be capable of similar feats from the three-dimensional perspective. Rudy Rucker illustrates this in his novel Spaceland, in which the protagonist encounters four-dimensional beings who demonstrate such powers.
Cross-sections
As a three-dimensional object passes through a two-dimensional plane, two-dimensional beings in this plane would only observe a cross-section of the three-dimensional object within this plane. For example, if a sphere passed through a sheet of paper, beings in the paper would see first a single point. A circle gradually grows larger, until it reaches the diameter of the sphere, and then gets smaller again, until it shrinks to a point and disappears. The 2D beings would not see a circle in the same way as three-dimensional beings do; rather, they only see a one-dimensional projection of the circle on their 1D "retina". Similarly, if a four-dimensional object passed through a three-dimensional (hyper) surface, one could observe a three-dimensional cross-section of the four-dimensional object. For example, a hypersphere would appear first as a point, then as a growing sphere (until it reaches the "hyperdiameter" of the hypersphere), with the sphere then shrinking to a single point and then disappearing. This means of visualizing aspects of the fourth dimension was used in the novel Flatland and also in several works of Charles Howard Hinton. And, in the same way, three-dimensional beings (such as humans with a 2D retina) can see all the sides and the insides of a 2D shape simultaneously, a 4D being could see all faces and the inside of a 3D shape at once with their 3D retina.
Projections
A useful application of dimensional analogy in visualizing higher dimensions is in projection. A projection is a way of representing an n-dimensional object in dimensions. For instance, computer screens are two-dimensional, and all the photographs of three-dimensional people, places, and things are represented in two dimensions by projecting the objects onto a flat surface. By doing this, the dimension orthogonal to the screen (depth) is removed and replaced with indirect information. The retina of the eye is also a two-dimensional array of receptors but the brain can perceive the nature of three-dimensional objects by inference from indirect information (such as shading, foreshortening, binocular vision, etc.). Artists often use perspective to give an illusion of three-dimensional depth to two-dimensional pictures. The shadow, cast by a fictitious grid model of a rotating tesseract on a plane surface, as shown in the figures, is also the result of projections.
Similarly, objects in the fourth dimension can be mathematically projected to the familiar three dimensions, where they can be more conveniently examined. In this case, the 'retina' of the four-dimensional eye is a three-dimensional array of receptors. A hypothetical being with such an eye would perceive the nature of four-dimensional objects by inferring four-dimensional depth from indirect information in the three-dimensional images in its retina.
The perspective projection of three-dimensional objects into the retina of the eye introduces artifacts such as foreshortening, which the brain interprets as depth in the third dimension. In the same way, perspective projection from four dimensions produces similar foreshortening effects. By applying dimensional analogy, one may infer four-dimensional "depth" from these effects.
As an illustration of this principle, the following sequence of images compares various views of the three-dimensional cube with analogous projections of the four-dimensional tesseract into three-dimensional space.
Shadows
A concept closely related to projection is the casting of shadows.
If a light is shone on a three-dimensional object, a two-dimensional shadow is cast. By dimensional analogy, light shone on a two-dimensional object in a two-dimensional world would cast a one-dimensional shadow, and light on a one-dimensional object in a one-dimensional world would cast a zero-dimensional shadow, that is, a point of non-light. Going the other way, one may infer that light shining on a four-dimensional object in a four-dimensional world would cast a three-dimensional shadow.
If the wireframe of a cube is lit from above, the resulting shadow on a flat two-dimensional surface is a square within a square with the corresponding corners connected. Similarly, if the wireframe of a tesseract were lit from "above" (in the fourth dimension), its shadow would be that of a three-dimensional cube within another three-dimensional cube suspended in midair (a "flat" surface from a four-dimensional perspective). (Note that, technically, the visual representation shown here is a two-dimensional image of the three-dimensional shadow of the four-dimensional wireframe figure.)
Bounding regions
The dimensional analogy also helps in inferring basic properties of objects in higher dimensions, such as the bounding region. For example, two-dimensional objects are bounded by one-dimensional boundaries: a square is bounded by four edges. Three-dimensional objects are bounded by two-dimensional surfaces: a cube is bounded by 6 square faces.
By applying dimensional analogy, one may infer that a four-dimensional cube, known as a tesseract, is bounded by three-dimensional volumes. And indeed, this is the case: mathematics shows that the tesseract is bounded by 8 cubes. Knowing this is key to understanding how to interpret a three-dimensional projection of the tesseract. The boundaries of the tesseract project to volumes in the image, not merely two-dimensional surfaces.
Hypervolume
The 4-volume or hypervolume in 4D can be calculated in closed form for simple geometrical figures, such as the tesseract (s4, for side length s) and the 4-ball ( for radius r).
Reasoning by analogy from familiar lower dimensions can be an excellent intuitive guide, but care must be exercised not to accept results that are not more rigorously tested. For example, consider the formulas for the area enclosed by a circle in two dimensions () and the volume enclosed by a sphere in three dimensions (). One might guess that the volume enclosed by the sphere in four-dimensional space is a rational multiple of , but the correct volume is . The volume of an n-ball in an arbitrary dimension n is computable from a recurrence relation connecting dimension to dimension .
In culture
In art
In literature
Science fiction texts often mention the concept of "dimension" when referring to parallel or alternate universes or other imagined planes of existence. This usage is derived from the idea that to travel to parallel/alternate universes/planes of existence one must travel in a direction/dimension besides the standard ones. In effect, the other universes/planes are just a small distance away from our own, but the distance is in a fourth (or higher) spatial (or non-spatial) dimension, not the standard ones.
One of the most heralded science fiction stories regarding true geometric dimensionality, and often recommended as a starting point for those just starting to investigate such matters, is the 1884 novella Flatland by Edwin A. Abbott. Isaac Asimov, in his foreword to the Signet Classics 1984 edition, described Flatland as "The best introduction one can find into the manner of perceiving dimensions."
The idea of other dimensions was incorporated into many early science fiction stories, appearing prominently, for example, in Miles J. Breuer's The Appendix and the Spectacles (1928) and Murray Leinster's The Fifth-Dimension Catapult (1931); and appeared irregularly in science fiction by the 1940s. Classic stories involving other dimensions include Robert A. Heinlein's —And He Built a Crooked House (1941), in which a California architect designs a house based on a three-dimensional projection of a tesseract; Alan E. Nourse's Tiger by the Tail and The Universe Between (both 1951); and The Ifth of Oofth (1957) by Walter Tevis. Another reference is Madeleine L'Engle's novel A Wrinkle In Time (1962), which uses the fifth dimension as a way of "tesseracting the universe" or "folding" space to move across it quickly. The fourth and fifth dimensions are also key components of the book The Boy Who Reversed Himself by William Sleator.
In philosophy
Immanuel Kant wrote in 1783: "That everywhere space (which is not itself the boundary of another space) has three dimensions and that space, in general, cannot have more dimensions is based on the proposition that not more than three lines can intersect at right angles in one point. This proposition cannot at all be shown from concepts, but rests immediately on intuition and indeed on pure intuition a priori because it is apodictically (demonstrably) certain."
"Space has Four Dimensions" is a short story published in 1846 by German philosopher and experimental psychologist Gustav Fechner under the pseudonym "Dr. Mises". The protagonist in the tale is a shadow who is aware of and able to communicate with other shadows, but who is trapped on a two-dimensional surface. According to Fechner, this "shadow-man" would conceive of the third dimension as being one of time. The story bears a strong similarity to the "Allegory of the Cave" presented in Plato's The Republic ( 380 BC).
Simon Newcomb wrote an article for the Bulletin of the American Mathematical Society in 1898 entitled "The Philosophy of Hyperspace". Linda Dalrymple Henderson coined the term "hyperspace philosophy", used to describe writing that uses higher dimensions to explore metaphysical themes, in her 1983 thesis about the fourth dimension in early-twentieth-century art. Examples of "hyperspace philosophers" include Charles Howard Hinton, the first writer, in 1888, to use the word "tesseract"; and the Russian esotericist P. D. Ouspensky.
| Mathematics | Four-dimensional space | null |
1365637 | https://en.wikipedia.org/wiki/Sports%20medicine | Sports medicine | Sports medicine is a branch of medicine that deals with physical fitness and the treatment and prevention of injuries related to sports and exercise. Although most sports teams have employed team physicians for many years, it is only since the late 20th century that sports medicine emerged as a distinct field of health care. In many countries, now over 50, sports medicine (or sport and exercise medicine) is a recognized medical specialty (with similar training and standards to other medical specialties or sub-specialties). In the majority of countries where sports medicine is recognized and practiced, it is a physician (non-surgical) specialty, but in some (such as the USA), it can equally be a surgical or non-surgical medical specialty, and also a specialty field within primary care. In other contexts, the field of sports medicine encompasses the scope of both medical specialists as well as allied health practitioners who work in the field of sport, such as physiotherapists, athletic trainers, podiatrists and exercise physiologists.
Scope
Sports medicine can refer to the specific medical specialty or subspecialty of several medical and research disciplines in sports. Sports medicine may be called Sport and Exercise medicine (SEM), which is now well established in many countries. It can broadly also refer to physicians, scientists, trainers, and other paramedical practitioners who work in a broad setting. Sports medicine specialists include a broad range of professions. All sports medicine specialists have one main goal in mind, and that is preventing future injuries and to improve the function of that area to return to everyday life. They work with all different types of people, and not just athletes. The various sports medicine experts often work together as a team to ensure the best recovery plan for the individual. Team members can include orthopedic surgeons, certified athletic trainers, sports physical therapists, physical medicine and rehabilitation specialists, and specialty SEM physicians.
Specializing in the treatment of athletes and other physically active individuals, SEM physicians have extensive education in musculoskeletal medicine. SEM doctors treat injuries such as muscle, ligament, tendon and bone problems, but may also treat chronic illnesses that can affect physical performance, such as asthma and diabetes. SEM doctors also advise on managing and preventing injuries.
European templates for SEM specialization generally recommend four years of experience in:
internal medicine with special emphasis on cardiology, emergency medicine and clinical nutrition
orthopedics and traumatology
physical and rehabilitation medicine
fellowship at a recognized sports medicine centre.
Related medical specialties
Exercise medicine
Podiatry
Sports cardiology
Emergency medicine
Lifestyle medicine
General practice
Adolescent medicine
Physical medicine and rehabilitation
Rheumatology
Orthopaedic sports medicine
Pediatrics
Occupational medicine
Establishment as a medical specialty
Historical roots of sports medicine
Although sports medicine was only established formally as a specialty in the 20th Century, the history of doctors having involvement in treating athletes goes back to ancient times in Greek, Roman and Egyptian societies.
Modern establishment of the specialty
Continental European countries were the first to establish medical groups with a focus on sport in the earliest part of the 20th Century. Possibly the earliest establishment of a society of Sports Medicine was the DGSP in Germany in 1912. The Italian version of this page Medicina dello sport states that Sports Medicine societies were first established in Switzerland (1922) followed by France (1929) and Italy (1929) (Italian Sports Medicine Federation). In Germany in the 1920s, an attempt was made to upskill thousands of doctors and other health professionals in sport and exercise medicine, without establishing it as a distinct specialty at that stage, but it failed due to lack of funding in the Depression. Sports medicine was established as a distinct specialty in Italy, the first country to do so, in 1958. The European Union of Medical Specialists has defined necessary training requirements for the establishment of the specialty of Sports Medicine in a given European country. In May 2024, the EU approved cross recognition of sports medicine qualifications between 11 different countries. It is a goal of the European Federation of Sports Medicine Associations to eventually establish Sports Medicine as a specialty in all European countries.
In Australia and New Zealand, Sport and Exercise Medicine (SEM) is a stand-alone medical specialty, with the Australasian College of Sport and Exercise Physicians being one of Australia's 15 recognized medical specialty Colleges. Australia, New Zealand and the UK have been cited as pioneer countries in the establishment of SEM as a stand-alone specialty.
The USA (and many other countries) follow the model of recognizing Sports Medicine as an official subspecialty of multiple other primary medical specialties. The most common primary specialties prior to a sports medicine subspecialty in the USA are family practice, orthopedics and physiatry.
Public health
SEM physicians are frequently involved in promoting the therapeutic benefits of physical activity, exercise and sport for the individuals and communities. SEM Physicians in the UK spend a period of their training in public health, and advise public health physicians on matters relating to physical activity promotion.
Common sports injuries
Common sports injuries that can result in seeing a sports medicine specialist are knee and shoulder injuries, fractures, ankle sprains, concussions, cartilage injuries, and more. A sports medicine specialist can also be seen for advice in other areas of health, like nutrition, exercise, supplements, and how to prevent injuries before they occur. A sports medicine specialist works to help make the performance of the athlete more advanced, as well as ensuring their safety while performing the activity. Sports injuries generally affect soft tissue or bones within the body and are commonly treated without surgery.
Treatment for sports injuries
Different types of sports injuries require different treatments and major injuries involve surgery, but most do not. Common treatments include medication, such as pain relievers or anti-inflammatory medication, icing, physical therapy, and/or immobilization of the injured area. Physical therapy is used to get the injured area back into regular movements and to reduce the discomfort of the affected area. PRICE is an acronym that is used for the common treatment of these injuries. It stands for protection, rest, ice, compression, and elevation.
Controversies in sports medicine
Concussion in sport
The management of concussion in sport has been extremely controversial over the past 20 years due to the discovery and reporting of Chronic traumatic encephalopathy as a disease that is common in ex-athletes, particularly footballers. Sporting codes have been accused of being complicit in understating the long-term damage caused by concussions by allowing too many head impacts to occur and for the players to be able to return to play too quickly after received concussions. A seminal series of consensus papers has been the international guidelines on the management of concussion in sport. These consensus statements have been seen on the positive side as being sports medicine leaders moving the management of concussion in a more conservative direction over time and encouraging a standard set of tests and assessments. On the negative side, they have been seen as conflicted and allowing return to play too rapidly.
Transgender people in sport
Whether male-to-female transgender athletes can safely and fairly participate in women's sport at the elite and community levels is a highly charged and controversial topic. The sports medicine world is not united in its views and although this debate well and truly involves medical input, it is as much a social controversy as it is a medical one.
Drugs in sport
Doping in sport has a long history with doctors in the sports medicine world being both heroes and villains on different occasions. The presence of trained sports medicine professionals at elite sporting events has been critical in the fight against doping, but sometimes doctors become the enablers of doping and are part of the scandal themselves.
Sports scandals involving medicine
Major scandals where doctors were prominent include:
Bloodgate
USA Gymnastics sex abuse scandal
Operación Puerto doping case
Biogenesis scandal
Eva Carneiro scandal, involving Chelsea F.C. and José Mourinho
Allied health team members
Different medical professionals for sports injuries require different forms of training, but for sports injuries, they mainly all work with the diagnosis and treatment of these injuries. All sports medicine professionals work with people of all age ranges, professional athletes, or even adolescents playing any sport. The main two allied health professions for sports injuries are athletic trainers (in the USA) and physical therapists (physiotherapists) in most other countries.
Athletic trainer
Athletic trainers are typically part of a sports medicine team in the US in particular, providing primary care, injury and illness prevention, wellness promotion, emergency care, therapeutic intervention and rehabilitation to injuries. When an athlete is injured, an athletic trainer is key to treatment and rehabilitation working closely with the athlete throughout rehabilitation. Athletic trainers are often the ones who assess the injury first and provide initial care.
Physiotherapist
Physiotherapists are a primary sports medicine team member in most countries of the world. Physiotherapists can specialize in many areas with sports physiotherapy as a major subspecialty. Physiotherapists are a main factor in the recovery stage of an injury as they set up an individualized recovery plan. Physiotherapy is underfunded within most health systems so that it is generally much more accessible in higher-income countries and, even within these countries, is much more accessible to higher-income earners. In countries like Denmark and Australia there are many more physiotherapists than in lower-income countries.
Podiatrist
Podiatrists treat issues related to the foot or ankle, which is a common area where athletes get injuries. They specialize in the diagnosis and treatment of foot-related issues by performing tests and referring physical therapists. Podiatrists can also perform surgeries or prescribe medication as forms of treatment.
Other practitioners
All of Exercise physiologists, Strength and conditioning coaches, personal trainers, Chiropractors, Osteopaths, Sports psychologists, Sekkotsu, and Sports nutritionists/dietitians can be part of the Sport and Exercise Medicine team.
Journals
| Biology and health sciences | Fields of medicine | Health |
1367595 | https://en.wikipedia.org/wiki/Elemental%20analysis | Elemental analysis | Elemental analysis is a process where a sample of some material (e.g., soil, waste or drinking water, bodily fluids, minerals, chemical compounds) is analyzed for its elemental and sometimes isotopic composition. Elemental analysis can be qualitative (determining what elements are present), and it can be quantitative (determining how much of each is present). Elemental analysis falls within the ambit of analytical chemistry, the instruments involved in deciphering the chemical nature of our world.
History
Antoine Lavoisier is regarded as the inventor of elemental analysis as a quantitative, experimental tool to assess the chemical composition of a compound. At the time, elemental analysis was based on the gravimetric determination of specific absorbent materials before and after selective adsorption of the combustion gases. Today fully automated systems based on thermal conductivity or infrared spectroscopy detection of the combustion gases, or other spectroscopic methods are used.
CHNX analysis
For organic chemists, elemental analysis or "EA" almost always refers to CHNX analysis—the determination of the mass fractions of carbon, hydrogen, nitrogen, and heteroatoms (X) (halogens, sulfur) of a sample. This information is important to help determine the structure of an unknown compound, as well as to help ascertain the structure and purity of a synthesized compound. In present-day organic chemistry, spectroscopic techniques (NMR, both 1H and 13C), mass spectrometry and chromatographic procedures have replaced EA as the primary technique for structural determination. However, it still gives very useful complementary information.
The most common form of elemental analysis, CHNS analysis, is accomplished by combustion analysis. Modern elemental analyzers are also capable of simultaneous determination of sulfur along with CHN in the same measurement run.
Quantitative analysis
Quantitative analysis determines the mass of each element or compound present. Other quantitative methods include gravimetry, optical atomic spectroscopy, and neutron activation analysis.
Gravimetry is where the sample is dissolved, the element of interest is precipitated and its mass measured, or the element of interest is volatilized, and the mass loss is measured.
Optical atomic spectroscopy includes flame atomic absorption, graphite furnace atomic absorption, and inductively coupled plasma atomic emission spectroscopy, which probe the outer electronic structure of atoms.
Neutron activation analysis involves the activation of a sample matrix through the process of neutron capture. The resulting radioactive target nuclei of the sample begin to decay, emitting gamma rays of specific energies that identify the radioisotopes present in the sample. The concentration of each analyte can be determined by comparison to an irradiated standard with known concentrations of each analyte.
Qualitative analysis
To qualitatively determine which elements exist in a sample, the methods are mass spectrometric atomic spectroscopy, such as inductively coupled plasma mass spectrometry, which probes the mass of atoms; other spectroscopy, which probes the inner electronic structure of atoms such as X-ray fluorescence, particle-induced X-ray emission, X-ray photoelectron spectroscopy, and Auger electron spectroscopy; and chemical methods such as the sodium fusion test and Schöniger oxidation.
Analysis of results
The analysis of results is performed by determining the ratio of elements from within the sample and working out a chemical formula that fits with those results. This process is useful as it helps determine if a sample sent is the desired compound and confirms the purity of a compound. The accepted deviation of elemental analysis results from the calculated is 0.3%.
| Physical sciences | Basics_2 | Chemistry |
1367789 | https://en.wikipedia.org/wiki/Monitor%20%28synchronization%29 | Monitor (synchronization) | In concurrent programming, a monitor is a synchronization construct that prevents threads from concurrently accessing a shared object's state and allows them to wait for the state to change. They provide a mechanism for threads to temporarily give up exclusive access in order to wait for some condition to be met, before regaining exclusive access and resuming their task. A monitor consists of a mutex (lock) and at least one condition variable. A condition variable is explicitly 'signalled' when the object's state is modified, temporarily passing the mutex to another thread 'waiting' on the condition variable.
Another definition of monitor is a thread-safe class, object, or module that wraps around a mutex in order to safely allow access to a method or variable by more than one thread. The defining characteristic of a monitor is that its methods are executed with mutual exclusion: At each point in time, at most one thread may be executing any of its methods. By using one or more condition variables it can also provide the ability for threads to wait on a certain condition (thus using the above definition of a "monitor"). For the rest of this article, this sense of "monitor" will be referred to as a "thread-safe object/class/module".
Monitors were invented by Per Brinch Hansen and C. A. R. Hoare, and were first implemented in Brinch Hansen's Concurrent Pascal language.
Mutual exclusion
While a thread is executing a method of a thread-safe object, it is said to occupy the object, by holding its mutex (lock). Thread-safe objects are implemented to enforce that at each point in time, at most one thread may occupy the object. The lock, which is initially unlocked, is locked at the start of each public method, and is unlocked at each return from each public method.
Upon calling one of the methods, a thread must wait until no other thread is executing any of the thread-safe object's methods before starting execution of its method. Note that without this mutual exclusion, two threads could cause money to be lost or gained for no reason. For example, two threads withdrawing 1000 from the account could both return true, while causing the balance to drop by only 1000, as follows: first, both threads fetch the current balance, find it greater than 1000, and subtract 1000 from it; then, both threads store the balance and return.
Condition variables
Problem statement
For many applications, mutual exclusion is not enough. Threads attempting an operation may need to wait until some condition holds true. A busy waiting loop
while not ( ) do skip
will not work, as mutual exclusion will prevent any other thread from entering the monitor to make the condition true. Other "solutions" exist such as having a loop that unlocks the monitor, waits a certain amount of time, locks the monitor and checks for the condition . Theoretically, it works and will not deadlock, but issues arise. It is hard to decide an appropriate amount of waiting time: too small and the thread will hog the CPU, too big and it will be apparently unresponsive. What is needed is a way to signal the thread when the condition is true (or could be true).
Case study: classic bounded producer/consumer problem
A classic concurrency problem is that of the bounded producer/consumer, in which there is a queue or ring buffer of tasks with a maximum size, with one or more threads being "producer" threads that add tasks to the queue, and one or more other threads being "consumer" threads that take tasks out of the queue. The queue is assumed to be non–thread-safe itself, and it can be empty, full, or between empty and full. Whenever the queue is full of tasks, then we need the producer threads to block until there is room from consumer threads dequeueing tasks. On the other hand, whenever the queue is empty, then we need the consumer threads to block until more tasks are available due to producer threads adding them.
As the queue is a concurrent object shared between threads, accesses to it must be made atomic, because the queue can be put into an inconsistent state during the course of the queue access that should never be exposed between threads. Thus, any code that accesses the queue constitutes a critical section that must be synchronized by mutual exclusion. If code and processor instructions in critical sections of code that access the queue could be interleaved by arbitrary context switches between threads on the same processor or by simultaneously-running threads on multiple processors, then there is a risk of exposing inconsistent state and causing race conditions.
Incorrect without synchronization
A naive approach is to design the code with busy-waiting and no synchronization, making the code subject to race conditions:
global RingBuffer queue; // A thread-unsafe ring-buffer of tasks.
// Method representing each producer thread's behavior:
public method producer() {
while (true) {
task myTask = ...; // Producer makes some new task to be added.
while (queue.isFull()) {} // Busy-wait until the queue is non-full.
queue.enqueue(myTask); // Add the task to the queue.
}
}
// Method representing each consumer thread's behavior:
public method consumer() {
while (true) {
while (queue.isEmpty()) {} // Busy-wait until the queue is non-empty.
myTask = queue.dequeue(); // Take a task off of the queue.
doStuff(myTask); // Go off and do something with the task.
}
}
This code has a serious problem in that accesses to the queue can be interrupted and interleaved with other threads' accesses to the queue. The queue.enqueue and queue.dequeue methods likely have instructions to update the queue's member variables such as its size, beginning and ending positions, assignment and allocation of queue elements, etc. In addition, the queue.isEmpty() and queue.isFull() methods read this shared state as well. If producer/consumer threads are allowed to be interleaved during the calls to enqueue/dequeue, then inconsistent state of the queue can be exposed leading to race conditions. In addition, if one consumer makes the queue empty in-between another consumer's exiting the busy-wait and calling "dequeue", then the second consumer will attempt to dequeue from an empty queue leading to an error. Likewise, if a producer makes the queue full in-between another producer's exiting the busy-wait and calling "enqueue", then the second producer will attempt to add to a full queue leading to an error.
Spin-waiting
One naive approach to achieve synchronization, as alluded to above, is to use "spin-waiting", in which a mutex is used to protect the critical sections of code and busy-waiting is still used, with the lock being acquired and released in between each busy-wait check.
global RingBuffer queue; // A thread-unsafe ring-buffer of tasks.
global Lock queueLock; // A mutex for the ring-buffer of tasks.
// Method representing each producer thread's behavior:
public method producer() {
while (true) {
task myTask = ...; // Producer makes some new task to be added.
queueLock.acquire(); // Acquire lock for initial busy-wait check.
while (queue.isFull()) { // Busy-wait until the queue is non-full.
queueLock.release();
// Drop the lock temporarily to allow a chance for other threads
// needing queueLock to run so that a consumer might take a task.
queueLock.acquire(); // Re-acquire the lock for the next call to "queue.isFull()".
}
queue.enqueue(myTask); // Add the task to the queue.
queueLock.release(); // Drop the queue lock until we need it again to add the next task.
}
}
// Method representing each consumer thread's behavior:
public method consumer() {
while (true) {
queueLock.acquire(); // Acquire lock for initial busy-wait check.
while (queue.isEmpty()) { // Busy-wait until the queue is non-empty.
queueLock.release();
// Drop the lock temporarily to allow a chance for other threads
// needing queueLock to run so that a producer might add a task.
queueLock.acquire(); // Re-acquire the lock for the next call to "queue.isEmpty()".
}
myTask = queue.dequeue(); // Take a task off of the queue.
queueLock.release(); // Drop the queue lock until we need it again to take off the next task.
doStuff(myTask); // Go off and do something with the task.
}
}
This method assures that an inconsistent state does not occur, but wastes CPU resources due to the unnecessary busy-waiting. Even if the queue is empty and producer threads have nothing to add for a long time, consumer threads are always busy-waiting unnecessarily. Likewise, even if consumers are blocked for a long time on processing their current tasks and the queue is full, producers are always busy-waiting. This is a wasteful mechanism. What is needed is a way to make producer threads block until the queue is non-full, and a way to make consumer threads block until the queue is non-empty.
(N.B.: Mutexes themselves can also be spin-locks which involve busy-waiting in order to get the lock, but in order to solve this problem of wasted CPU resources, we assume that queueLock is not a spin-lock and properly uses a blocking lock queue itself.)
Condition variables
The solution is to use condition variables. Conceptually, a condition variable is a queue of threads, associated with a mutex, on which a thread may wait for some condition to become true. Thus each condition variable is associated with an assertion . While a thread is waiting on a condition variable, that thread is not considered to occupy the monitor, and so other threads may enter the monitor to change the monitor's state. In most types of monitors, these other threads may signal the condition variable to indicate that assertion is true in the current state.
Thus there are three main operations on condition variables:
wait c, m, where c is a condition variable and m is a mutex (lock) associated with the monitor. This operation is called by a thread that needs to wait until the assertion is true before proceeding. While the thread is waiting, it does not occupy the monitor. The function, and fundamental contract, of the "wait" operation, is to do the following steps:
Atomically:
Once this thread is subsequently notified/signaled (see below) and resumed, then automatically re-acquire the mutex m.
Steps 1a and 1b can occur in either order, with 1c usually occurring after them. While the thread is sleeping and in c's wait-queue, the next program counter to be executed is at step 2, in the middle of the "wait" function/subroutine. Thus, the thread sleeps and later wakes up in the middle of the "wait" operation.
The atomicity of the operations within step 1 is important to avoid race conditions that would be caused by a preemptive thread switch in-between them. One failure mode that could occur if these were not atomic is a missed wakeup, in which the thread could be on c's sleep-queue and have released the mutex, but a preemptive thread switch occurred before the thread went to sleep, and another thread called a signal operation (see below) on c moving the first thread back out of c's queue. As soon as the first thread in question is switched back to, its program counter will be at step 1c, and it will sleep and be unable to be woken up again, violating the invariant that it should have been on c's sleep-queue when it slept. Other race conditions depend on the ordering of steps 1a and 1b, and depend on where a context switch occurs.
signal c, also known as notify c, is called by a thread to indicate that the assertion is true. Depending on the type and implementation of the monitor, this moves one or more threads from c's sleep-queue to the "ready queue", or another queue for it to be executed. It is usually considered a best practice to perform the "signal" operation before releasing mutex m that is associated with c, but as long as the code is properly designed for concurrency and depending on the threading implementation, it is often also acceptable to release the lock before signalling. Depending on the threading implementation, the ordering of this can have scheduling-priority ramifications. (Some authors instead advocate a preference for releasing the lock before signalling.) A threading implementation should document any special constraints on this ordering.
broadcast c, also known as notifyAll c, is a similar operation that wakes up all threads in c's wait-queue. This empties the wait-queue. Generally, when more than one predicate condition is associated with the same condition variable, the application will require broadcast instead of signal because a thread waiting for the wrong condition might be woken up and then immediately go back to sleep without waking up a thread waiting for the correct condition that just became true. Otherwise, if the predicate condition is one-to-one with the condition variable associated with it, then signal may be more efficient than broadcast.
As a design rule, multiple condition variables can be associated with the same mutex, but not vice versa. (This is a one-to-many correspondence.) This is because the predicate is the same for all threads using the monitor and must be protected with mutual exclusion from all other threads that might cause the condition to be changed or that might read it while the thread in question causes it to be changed, but there may be different threads that want to wait for a different condition on the same variable requiring the same mutex to be used. In the producer-consumer example described above, the queue must be protected by a unique mutex object, m. The "producer" threads will want to wait on a monitor using lock m and a condition variable which blocks until the queue is non-full. The "consumer" threads will want to wait on a different monitor using the same mutex m but a different condition variable which blocks until the queue is non-empty. It would (usually) never make sense to have different mutexes for the same condition variable, but this classic example shows why it often certainly makes sense to have multiple condition variables using the same mutex. A mutex used by one or more condition variables (one or more monitors) may also be shared with code that does not use condition variables (and which simply acquires/releases it without any wait/signal operations), if those critical sections do not happen to require waiting for a certain condition on the concurrent data.
Monitor usage
The proper basic usage of a monitor is:
acquire(m); // Acquire this monitor's lock.
while (!p) { // While the condition/predicate/assertion that we are waiting for is not true...
wait(m, cv); // Wait on this monitor's lock and condition variable.
}
// ... Critical section of code goes here ...
signal(cv2); // Or: broadcast(cv2);
// cv2 might be the same as cv or different.
release(m); // Release this monitor's lock.
The following is the same pseudocode but with more verbose comments to better explain what is going on:
// ... (previous code)
// About to enter the monitor.
// Acquire the advisory mutex (lock) associated with the concurrent
// data that is shared between threads,
// to ensure that no two threads can be preemptively interleaved or
// run simultaneously on different cores while executing in critical
// sections that read or write this same concurrent data. If another
// thread is holding this mutex, then this thread will be put to sleep
// (blocked) and placed on m's sleep queue. (Mutex "m" shall not be
// a spin-lock.)
acquire(m);
// Now, we are holding the lock and can check the condition for the
// first time.
// The first time we execute the while loop condition after the above
// "acquire", we are asking, "Does the condition/predicate/assertion
// we are waiting for happen to already be true?"
while (!p()) // "p" is any expression (e.g. variable or
// function-call) that checks the condition and
// evaluates to boolean. This itself is a critical
// section, so you *MUST* be holding the lock when
// executing this "while" loop condition!
// If this is not the first time the "while" condition is being checked,
// then we are asking the question, "Now that another thread using this
// monitor has notified me and woken me up and I have been context-switched
// back to, did the condition/predicate/assertion we are waiting on stay
// true between the time that I was woken up and the time that I re-acquired
// the lock inside the "wait" call in the last iteration of this loop, or
// did some other thread cause the condition to become false again in the
// meantime thus making this a spurious wakeup?
{
// If this is the first iteration of the loop, then the answer is
// "no" -- the condition is not ready yet. Otherwise, the answer is:
// the latter. This was a spurious wakeup, some other thread occurred
// first and caused the condition to become false again, and we must
// wait again.
wait(m, cv);
// Temporarily prevent any other thread on any core from doing
// operations on m or cv.
// release(m) // Atomically release lock "m" so other
// // code using this concurrent data
// // can operate, move this thread to cv's
// // wait-queue so that it will be notified
// // sometime when the condition becomes
// // true, and sleep this thread. Re-enable
// // other threads and cores to do
// // operations on m and cv.
//
// Context switch occurs on this core.
//
// At some future time, the condition we are waiting for becomes
// true, and another thread using this monitor (m, cv) does either
// a signal that happens to wake this thread up, or a
// broadcast that wakes us up, meaning that we have been taken out
// of cv's wait-queue.
//
// During this time, other threads may cause the condition to
// become false again, or the condition may toggle one or more
// times, or it may happen to stay true.
//
// This thread is switched back to on some core.
//
// acquire(m) // Lock "m" is re-acquired.
// End this loop iteration and re-check the "while" loop condition to make
// sure the predicate is still true.
}
// The condition we are waiting for is true!
// We are still holding the lock, either from before entering the monitor or from
// the last execution of "wait".
// Critical section of code goes here, which has a precondition that our predicate
// must be true.
// This code might make cv's condition false, and/or make other condition variables'
// predicates true.
// Call signal or broadcast, depending on which condition variables'
// predicates (who share mutex m) have been made true or may have been made true,
// and the monitor semantic type being used.
for (cv_x in cvs_to_signal) {
signal(cv_x); // Or: broadcast(cv_x);
}
// One or more threads have been woken up but will block as soon as they try
// to acquire m.
// Release the mutex so that notified thread(s) and others can enter their critical
// sections.
release(m);
Solving the bounded producer/consumer problem
Having introduced the usage of condition variables, let us use it to revisit and solve the classic bounded producer/consumer problem. The classic solution is to use two monitors, comprising two condition variables sharing one lock on the queue:
global volatile RingBuffer queue; // A thread-unsafe ring-buffer of tasks.
global Lock queueLock; // A mutex for the ring-buffer of tasks. (Not a spin-lock.)
global CV queueEmptyCV; // A condition variable for consumer threads waiting for the queue to
// become non-empty. Its associated lock is "queueLock".
global CV queueFullCV; // A condition variable for producer threads waiting for the queue to
// become non-full. Its associated lock is also "queueLock".
// Method representing each producer thread's behavior:
public method producer() {
while (true) {
// Producer makes some new task to be added.
task myTask = ...;
// Acquire "queueLock" for the initial predicate check.
queueLock.acquire();
// Critical section that checks if the queue is non-full.
while (queue.isFull()) {
// Release "queueLock", enqueue this thread onto "queueFullCV" and sleep this thread.
wait(queueLock, queueFullCV);
// When this thread is awoken, re-acquire "queueLock" for the next predicate check.
}
// Critical section that adds the task to the queue (note that we are holding "queueLock").
queue.enqueue(myTask);
// Wake up one or all consumer threads that are waiting for the queue to be non-empty
// now that it is guaranteed, so that a consumer thread will take the task.
signal(queueEmptyCV); // Or: broadcast(queueEmptyCV);
// End of critical sections.
// Release "queueLock" until we need it again to add the next task.
queueLock.release();
}
}
// Method representing each consumer thread's behavior:
public method consumer() {
while (true) {
// Acquire "queueLock" for the initial predicate check.
queueLock.acquire();
// Critical section that checks if the queue is non-empty.
while (queue.isEmpty()) {
// Release "queueLock", enqueue this thread onto "queueEmptyCV" and sleep this thread.
wait(queueLock, queueEmptyCV);
// When this thread is awoken, re-acquire "queueLock" for the next predicate check.
}
// Critical section that takes a task off of the queue (note that we are holding "queueLock").
myTask = queue.dequeue();
// Wake up one or all producer threads that are waiting for the queue to be non-full
// now that it is guaranteed, so that a producer thread will add a task.
signal(queueFullCV); // Or: broadcast(queueFullCV);
// End of critical sections.
// Release "queueLock" until we need it again to take the next task.
queueLock.release();
// Go off and do something with the task.
doStuff(myTask);
}
}
This ensures concurrency between the producer and consumer threads sharing the task queue, and blocks the threads that have nothing to do rather than busy-waiting as shown in the aforementioned approach using spin-locks.
A variant of this solution could use a single condition variable for both producers and consumers, perhaps named "queueFullOrEmptyCV" or "queueSizeChangedCV". In this case, more than one condition is associated with the condition variable, such that the condition variable represents a weaker condition than the conditions being checked by individual threads. The condition variable represents threads that are waiting for the queue to be non-full and ones waiting for it to be non-empty. However, doing this would require using broadcast in all the threads using the condition variable and cannot use a regular signal. This is because the regular signal might wake up a thread of the wrong type whose condition has not yet been met, and that thread would go back to sleep without a thread of the correct type getting signaled. For example, a producer might make the queue full and wake up another producer instead of a consumer, and the woken producer would go back to sleep. In the complementary case, a consumer might make the queue empty and wake up another consumer instead of a producer, and the consumer would go back to sleep. Using broadcast ensures that some thread of the right type will proceed as expected by the problem statement.
Here is the variant using only one condition variable and broadcast:
global volatile RingBuffer queue; // A thread-unsafe ring-buffer of tasks.
global Lock queueLock; // A mutex for the ring-buffer of tasks. (Not a spin-lock.)
global CV queueFullOrEmptyCV; // A single condition variable for when the queue is not ready for any thread
// i.e. for producer threads waiting for the queue to become non-full
// and consumer threads waiting for the queue to become non-empty.
// Its associated lock is "queueLock".
// Not safe to use regular "signal" because it is associated with
// multiple predicate conditions (assertions).
// Method representing each producer thread's behavior:
public method producer() {
while (true) {
// Producer makes some new task to be added.
task myTask = ...;
// Acquire "queueLock" for the initial predicate check.
queueLock.acquire();
// Critical section that checks if the queue is non-full.
while (queue.isFull()) {
// Release "queueLock", enqueue this thread onto "queueFullOrEmptyCV" and sleep this thread.
wait(queueLock, queueFullOrEmptyCV);
// When this thread is awoken, re-acquire "queueLock" for the next predicate check.
}
// Critical section that adds the task to the queue (note that we are holding "queueLock").
queue.enqueue(myTask);
// Wake up all producer and consumer threads that are waiting for the queue to be respectively
// non-full and non-empty now that the latter is guaranteed, so that a consumer thread will take the task.
broadcast(queueFullOrEmptyCV); // Do not use "signal" (as it might wake up another producer thread only).
// End of critical sections.
// Release "queueLock" until we need it again to add the next task.
queueLock.release();
}
}
// Method representing each consumer thread's behavior:
public method consumer() {
while (true) {
// Acquire "queueLock" for the initial predicate check.
queueLock.acquire();
// Critical section that checks if the queue is non-empty.
while (queue.isEmpty()) {
// Release "queueLock", enqueue this thread onto "queueFullOrEmptyCV" and sleep this thread.
wait(queueLock, queueFullOrEmptyCV);
// When this thread is awoken, re-acquire "queueLock" for the next predicate check.
}
// Critical section that takes a task off of the queue (note that we are holding "queueLock").
myTask = queue.dequeue();
// Wake up all producer and consumer threads that are waiting for the queue to be respectively
// non-full and non-empty now that the former is guaranteed, so that a producer thread will add a task.
broadcast(queueFullOrEmptyCV); // Do not use "signal" (as it might wake up another consumer thread only).
// End of critical sections.
// Release "queueLock" until we need it again to take the next task.
queueLock.release();
// Go off and do something with the task.
doStuff(myTask);
}
}
Synchronization primitives
Monitors are implemented using an atomic read-modify-write primitive and a waiting primitive. The read-modify-write primitive (usually test-and-set or compare-and-swap) is usually in the form of a memory-locking instruction provided by the ISA, but can also be composed of non-locking instructions on single-processor devices when interrupts are disabled. The waiting primitive can be a busy-wait loop or an OS-provided primitive that prevents the thread from being scheduled until it is ready to proceed.
Here is an example pseudocode implementation of parts of a threading system and mutexes and Mesa-style condition variables, using test-and-set and a first-come, first-served policy:
Sample Mesa-monitor implementation with Test-and-Set
// Basic parts of threading system:
// Assume "ThreadQueue" supports random access.
public volatile ThreadQueue readyQueue; // Thread-unsafe queue of ready threads. Elements are (Thread*).
public volatile global Thread* currentThread; // Assume this variable is per-core. (Others are shared.)
// Implements a spin-lock on just the synchronized state of the threading system itself.
// This is used with test-and-set as the synchronization primitive.
public volatile global bool threadingSystemBusy = false;
// Context-switch interrupt service routine (ISR):
// On the current CPU core, preemptively switch to another thread.
public method contextSwitchISR() {
if (testAndSet(threadingSystemBusy)) {
return; // Can't switch context right now.
}
// Ensure this interrupt can't happen again which would foul up the context switch:
systemCall_disableInterrupts();
// Get all of the registers of the currently-running process.
// For Program Counter (PC), we will need the instruction location of
// the "resume" label below. Getting the register values is platform-dependent and may involve
// reading the current stack frame, JMP/CALL instructions, etc. (The details are beyond this scope.)
currentThread->registers = getAllRegisters(); // Store the registers in the "currentThread" object in memory.
currentThread->registers.PC = resume; // Set the next PC to the "resume" label below in this method.
readyQueue.enqueue(currentThread); // Put this thread back onto the ready queue for later execution.
Thread* otherThread = readyQueue.dequeue(); // Remove and get the next thread to run from the ready queue.
currentThread = otherThread; // Replace the global current-thread pointer value so it is ready for the next thread.
// Restore the registers from currentThread/otherThread, including a jump to the stored PC of the other thread
// (at "resume" below). Again, the details of how this is done are beyond this scope.
restoreRegisters(otherThread.registers);
// *** Now running "otherThread" (which is now "currentThread")! The original thread is now "sleeping". ***
resume: // This is where another contextSwitch() call needs to set PC to when switching context back here.
// Return to where otherThread left off.
threadingSystemBusy = false; // Must be an atomic assignment.
systemCall_enableInterrupts(); // Turn pre-emptive switching back on on this core.
}
// Thread sleep method:
// On current CPU core, a synchronous context switch to another thread without putting
// the current thread on the ready queue.
// Must be holding "threadingSystemBusy" and disabled interrupts so that this method
// doesn't get interrupted by the thread-switching timer which would call contextSwitchISR().
// After returning from this method, must clear "threadingSystemBusy".
public method threadSleep() {
// Get all of the registers of the currently-running process.
// For Program Counter (PC), we will need the instruction location of
// the "resume" label below. Getting the register values is platform-dependent and may involve
// reading the current stack frame, JMP/CALL instructions, etc. (The details are beyond this scope.)
currentThread->registers = getAllRegisters(); // Store the registers in the "currentThread" object in memory.
currentThread->registers.PC = resume; // Set the next PC to the "resume" label below in this method.
// Unlike contextSwitchISR(), we will not place currentThread back into readyQueue.
// Instead, it has already been placed onto a mutex's or condition variable's queue.
Thread* otherThread = readyQueue.dequeue(); // Remove and get the next thread to run from the ready queue.
currentThread = otherThread; // Replace the global current-thread pointer value so it is ready for the next thread.
// Restore the registers from currentThread/otherThread, including a jump to the stored PC of the other thread
// (at "resume" below). Again, the details of how this is done are beyond this scope.
restoreRegisters(otherThread.registers);
// *** Now running "otherThread" (which is now "currentThread")! The original thread is now "sleeping". ***
resume: // This is where another contextSwitch() call needs to set PC to when switching context back here.
// Return to where otherThread left off.
}
public method wait(Mutex m, ConditionVariable c) {
// Internal spin-lock while other threads on any core are accessing this object's
// "held" and "threadQueue", or "readyQueue".
while (testAndSet(threadingSystemBusy)) {}
// N.B.: "threadingSystemBusy" is now true.
// System call to disable interrupts on this core so that threadSleep() doesn't get interrupted by
// the thread-switching timer on this core which would call contextSwitchISR().
// Done outside threadSleep() for more efficiency so that this thread will be sleeped
// right after going on the condition-variable queue.
systemCall_disableInterrupts();
assert m.held; // (Specifically, this thread must be the one holding it.)
m.release();
c.waitingThreads.enqueue(currentThread);
threadSleep();
// Thread sleeps ... Thread gets woken up from a signal/broadcast.
threadingSystemBusy = false; // Must be an atomic assignment.
systemCall_enableInterrupts(); // Turn pre-emptive switching back on on this core.
// Mesa style:
// Context switches may now occur here, making the client caller's predicate false.
m.acquire();
}
public method signal(ConditionVariable c) {
// Internal spin-lock while other threads on any core are accessing this object's
// "held" and "threadQueue", or "readyQueue".
while (testAndSet(threadingSystemBusy)) {}
// N.B.: "threadingSystemBusy" is now true.
// System call to disable interrupts on this core so that threadSleep() doesn't get interrupted by
// the thread-switching timer on this core which would call contextSwitchISR().
// Done outside threadSleep() for more efficiency so that this thread will be sleeped
// right after going on the condition-variable queue.
systemCall_disableInterrupts();
if (!c.waitingThreads.isEmpty()) {
wokenThread = c.waitingThreads.dequeue();
readyQueue.enqueue(wokenThread);
}
threadingSystemBusy = false; // Must be an atomic assignment.
systemCall_enableInterrupts(); // Turn pre-emptive switching back on on this core.
// Mesa style:
// The woken thread is not given any priority.
}
public method broadcast(ConditionVariable c) {
// Internal spin-lock while other threads on any core are accessing this object's
// "held" and "threadQueue", or "readyQueue".
while (testAndSet(threadingSystemBusy)) {}
// N.B.: "threadingSystemBusy" is now true.
// System call to disable interrupts on this core so that threadSleep() doesn't get interrupted by
// the thread-switching timer on this core which would call contextSwitchISR().
// Done outside threadSleep() for more efficiency so that this thread will be sleeped
// right after going on the condition-variable queue.
systemCall_disableInterrupts();
while (!c.waitingThreads.isEmpty()) {
wokenThread = c.waitingThreads.dequeue();
readyQueue.enqueue(wokenThread);
}
threadingSystemBusy = false; // Must be an atomic assignment.
systemCall_enableInterrupts(); // Turn pre-emptive switching back on on this core.
// Mesa style:
// The woken threads are not given any priority.
}
class Mutex {
protected volatile bool held = false;
private volatile ThreadQueue blockingThreads; // Thread-unsafe queue of blocked threads. Elements are (Thread*).
public method acquire() {
// Internal spin-lock while other threads on any core are accessing this object's
// "held" and "threadQueue", or "readyQueue".
while (testAndSet(threadingSystemBusy)) {}
// N.B.: "threadingSystemBusy" is now true.
// System call to disable interrupts on this core so that threadSleep() doesn't get interrupted by
// the thread-switching timer on this core which would call contextSwitchISR().
// Done outside threadSleep() for more efficiency so that this thread will be sleeped
// right after going on the lock queue.
systemCall_disableInterrupts();
assert !blockingThreads.contains(currentThread);
if (held) {
// Put "currentThread" on this lock's queue so that it will be
// considered "sleeping" on this lock.
// Note that "currentThread" still needs to be handled by threadSleep().
readyQueue.remove(currentThread);
blockingThreads.enqueue(currentThread);
threadSleep();
// Now we are woken up, which must be because "held" became false.
assert !held;
assert !blockingThreads.contains(currentThread);
}
held = true;
threadingSystemBusy = false; // Must be an atomic assignment.
systemCall_enableInterrupts(); // Turn pre-emptive switching back on on this core.
}
public method release() {
// Internal spin-lock while other threads on any core are accessing this object's
// "held" and "threadQueue", or "readyQueue".
while (testAndSet(threadingSystemBusy)) {}
// N.B.: "threadingSystemBusy" is now true.
// System call to disable interrupts on this core for efficiency.
systemCall_disableInterrupts();
assert held; // (Release should only be performed while the lock is held.)
held = false;
if (!blockingThreads.isEmpty()) {
Thread* unblockedThread = blockingThreads.dequeue();
readyQueue.enqueue(unblockedThread);
}
threadingSystemBusy = false; // Must be an atomic assignment.
systemCall_enableInterrupts(); // Turn pre-emptive switching back on on this core.
}
}
struct ConditionVariable {
volatile ThreadQueue waitingThreads;
}
Blocking condition variables
The original proposals by C. A. R. Hoare and Per Brinch Hansen were for blocking condition variables. With a blocking condition variable, the signaling thread must wait outside the monitor (at least) until the signaled thread relinquishes occupancy of the monitor by either returning or by again waiting on a condition variable. Monitors using blocking condition variables are often called Hoare-style monitors or signal-and-urgent-wait monitors.
We assume there are two queues of threads associated with each monitor object
e is the entrance queue
s is a queue of threads that have signaled.
In addition we assume that for each condition variable , there is a queue
.q, which is a queue for threads waiting on condition variable
All queues are typically guaranteed to be fair and, in some implementations, may be guaranteed to be first in first out.
The implementation of each operation is as follows. (We assume that each operation runs in mutual exclusion to the others; thus restarted threads do not begin executing until the operation is complete.)
enter the monitor:
enter the method
if the monitor is locked
add this thread to e
block this thread
else
lock the monitor
leave the monitor:
schedule
return from the method
wait :
add this thread to .q
schedule
block this thread
signal :
if there is a thread waiting on .q
select and remove one such thread t from .q
(t is called "the signaled thread")
add this thread to s
restart t
(so t will occupy the monitor next)
block this thread
schedule:
if there is a thread on s
select and remove one thread from s and restart it
(this thread will occupy the monitor next)
else if there is a thread on e
select and remove one thread from e and restart it
(this thread will occupy the monitor next)
else
unlock the monitor
(the monitor will become unoccupied)
The schedule routine selects the next thread to occupy the monitor
or, in the absence of any candidate threads, unlocks the monitor.
The resulting signaling discipline is known as "signal and urgent wait," as the signaler must wait, but is given priority over threads on the entrance queue. An alternative is "signal and wait," in which there is no s queue and signaler waits on the e queue instead.
Some implementations provide a signal and return operation that combines signaling with returning from a procedure.
signal and return:
if there is a thread waiting on .q
select and remove one such thread t from .q
(t is called "the signaled thread")
restart t
(so t will occupy the monitor next)
else
schedule
return from the method
In either case ("signal and urgent wait" or "signal and wait"), when a condition variable is signaled and there is at least one thread waiting on the condition variable, the signaling thread hands occupancy over to the signaled thread seamlessly, so that no other thread can gain occupancy in between. If is true at the start of each signal operation, it will be true at the end of each wait operation. This is summarized by the following contracts. In these contracts, is the monitor's invariant.
enter the monitor:
postcondition
leave the monitor:
precondition
wait :
precondition
modifies the state of the monitor
postcondition and
signal :
precondition and
modifies the state of the monitor
postcondition
signal and return:
precondition and
In these contracts, it is assumed that and do not depend on the
contents or lengths of any queues.
(When the condition variable can be queried as to the number of threads waiting on its queue, more sophisticated contracts can be given. For example, a useful pair of contracts, allowing occupancy to be passed without establishing the invariant, is:
wait :
precondition
modifies the state of the monitor
postcondition
signal
precondition (not empty() and ) or (empty() and )
modifies the state of the monitor
postcondition
(See Howard and Buhr et al. for more.)
It is important to note here that the assertion is entirely up to the programmer; he or she simply needs to be consistent about what it is.
We conclude this section with an example of a thread-safe class using a blocking monitor that implements a bounded, thread-safe stack.
monitor class SharedStack {
private const capacity := 10
private int[capacity] A
private int size := 0
invariant 0 <= size and size <= capacity
private BlockingCondition theStackIsNotEmpty /* associated with 0 < size and size <= capacity */
private BlockingCondition theStackIsNotFull /* associated with 0 <= size and size < capacity */
public method push(int value)
{
if size = capacity then wait theStackIsNotFull
assert 0 <= size and size < capacity
A[size] := value ; size := size + 1
assert 0 < size and size <= capacity
signal theStackIsNotEmpty and return
}
public method int pop()
{
if size = 0 then wait theStackIsNotEmpty
assert 0 < size and size <= capacity
size := size - 1 ;
assert 0 <= size and size < capacity
signal theStackIsNotFull and return A[size]
}
}
Note that, in this example, the thread-safe stack is internally providing a mutex, which, as in the earlier producer/consumer example, is shared by both condition variables, which are checking different conditions on the same concurrent data. The only difference is that the producer/consumer example assumed a regular non-thread-safe queue and was using a standalone mutex and condition variables, without these details of the monitor abstracted away as is the case here. In this example, when the "wait" operation is called, it must somehow be supplied with the thread-safe stack's mutex, such as if the "wait" operation is an integrated part of the "monitor class". Aside from this kind of abstracted functionality, when a "raw" monitor is used, it will always have to include a mutex and a condition variable, with a unique mutex for each condition variable.
Nonblocking condition variables
With nonblocking condition variables (also called "Mesa style" condition variables or "signal and continue" condition variables), signaling does not cause the signaling thread to lose occupancy of the monitor. Instead the signaled threads are moved to the e queue. There is no need for the s queue.
With nonblocking condition variables, the signal operation is often called notify — a terminology we will follow here. It is also common to provide a notify all operation that moves all threads waiting on a condition variable to the e queue.
The meaning of various operations are given here. (We assume that each operation runs in mutual exclusion to the others; thus restarted threads do not begin executing until the operation is complete.)
enter the monitor:
enter the method
if the monitor is locked
add this thread to e
block this thread
else
lock the monitor
leave the monitor:
schedule
return from the method
wait :
add this thread to .q
schedule
block this thread
notify :
if there is a thread waiting on .q
select and remove one thread t from .q
(t is called "the notified thread")
move t to e
notify all :
move all threads waiting on .q to e
schedule :
if there is a thread on e
select and remove one thread from e and restart it
else
unlock the monitor
As a variation on this scheme, the notified thread may be moved to a queue called w, which has priority over e. See Howard and Buhr et al. for further discussion.
It is possible to associate an assertion with each condition variable such that is sure to be true upon return from wait . However, one must
ensure that is preserved from the time the notifying thread gives up occupancy until the notified thread is selected to re-enter the monitor. Between these times there could be activity by other occupants. Thus it is common for to simply be true.
For this reason, it is usually necessary to enclose each wait operation in a loop like this
while not ( ) do
wait c
where is some condition stronger than . The operations notify and notify all are treated as "hints" that may be true for some waiting thread.
Every iteration of such a loop past the first represents a lost notification; thus with nonblocking monitors, one must be careful to ensure that too many notifications cannot be lost.
As an example of "hinting," consider a bank account in which a withdrawing thread will wait until the account has sufficient funds before proceeding
monitor class Account {
private int balance := 0
invariant balance >= 0
private NonblockingCondition balanceMayBeBigEnough
public method withdraw(int amount)
precondition amount >= 0
{
while balance < amount do wait balanceMayBeBigEnough
assert balance >= amount
balance := balance - amount
}
public method deposit(int amount)
precondition amount >= 0
{
balance := balance + amount
notify all balanceMayBeBigEnough
}
}
In this example, the condition being waited for is a function of the amount to be withdrawn, so it is impossible for a depositing thread to know that it made such a condition true. It makes sense in this case to allow each waiting thread into the monitor (one at a time) to check if its assertion is true.
Implicit condition variable monitors
In the Java language, each object may be used as a monitor. Methods requiring mutual exclusion must be explicitly marked with the synchronized keyword. Blocks of code may also be marked by synchronized.
Rather than having explicit condition variables, each monitor (i.e., object) is equipped with a single wait queue in addition to its entrance queue. All waiting is done on this single wait queue and all notify and notifyAll operations apply to this queue. This approach has been adopted in other languages, for example C#.
Implicit signaling
Another approach to signaling is to omit the signal operation. Whenever a thread leaves the monitor (by returning or waiting), the assertions of all waiting threads are evaluated until one is found to be true. In such a system, condition variables are not needed, but the assertions must be explicitly coded. The contract for wait is
wait :
precondition
modifies the state of the monitor
postcondition and
History
Brinch Hansen and Hoare developed the monitor concept in the early 1970s, based on earlier ideas of their own and of Edsger Dijkstra. Brinch Hansen published the first monitor notation, adopting the class concept of Simula 67, and invented a queueing mechanism. Hoare refined the rules of process resumption. Brinch Hansen created the first implementation of monitors, in Concurrent Pascal. Hoare demonstrated their equivalence to semaphores.
Monitors (and Concurrent Pascal) were soon used to structure process synchronization in the Solo operating system.
Programming languages that have supported monitors include:
Ada since Ada 95 (as protected objects)
C# (and other languages that use the .NET Framework)
Concurrent Euclid
Concurrent Pascal
D
Delphi (Delphi 2009 and above, via TObject.Monitor)
Java (via the wait and notify methods)
Go
Mesa
Modula-3
Python (via threading.Condition object)
Ruby
Squeak Smalltalk
Turing, Turing+, and Object-Oriented Turing
μC++
Visual Prolog
A number of libraries have been written that allow monitors to be constructed in languages that do not support them natively. When library calls are used, it is up to the programmer to explicitly mark the start and end of code executed with mutual exclusion. Pthreads is one such library.
| Technology | Computer science | null |
1368466 | https://en.wikipedia.org/wiki/Bacterial%20capsule | Bacterial capsule | The bacterial capsule is a large structure common to many bacteria. It is a polysaccharide layer that lies outside the cell envelope, and is thus deemed part of the outer envelope of a bacterial cell. It is a well-organized layer, not easily washed off, and it can be the cause of various diseases.
The capsule—which can be found in both gram negative and gram-positive bacteria—is different from the second lipid membrane – bacterial outer membrane, which contains lipopolysaccharides and lipoproteins and is found only in gram-negative bacteria.
When the amorphous viscid secretion (that makes up the capsule) diffuses into the surrounding medium and remains as a loose undemarcated secretion, it is known as a slime layer. Capsule and slime layer are sometimes summarized under the term glycocalyx.
Composition
Most bacterial capsules are composed of polysaccharide, but some species use other materials, such as poly-D-glutamic acid in Bacillus anthracis. Because most capsules are so tightly packed, they are difficult to stain because most standard stains cannot penetrate the capsule. To visualize encapsulated bacteria using a microscope, a sample is treated with a dark stain, such as India ink. The structure of the capsule prevents the stain from penetrating the cell. When viewed,JB bacterial capsules appear as a bright halo around the cell on a dark background.
Function
The bacterial capsule serves as a shield, giving protection from toxins, and from drying out. Capsules allow adhesion to surfaces and help enable the bacteria to evade the host immune system. The water content in the capsule gives the protection against drying out.
The capsule is considered a virulence factor because it enhances the ability of bacteria to cause disease (e.g. prevents phagocytosis). The capsule can protect cells from engulfment by eukaryotic cells, such as macrophages. A capsule-specific antibody may be required for phagocytosis to occur. They also exclude bacterial viruses and most hydrophobic toxic materials such as detergents. Immunity to one capsule type does not result in immunity to the other types. Capsules also help cells adhere to surfaces. As a group where the capsule is present they are known as polysaccharide encapsulated bacteria or encapsulated bacteria.
Diversity
The capsule is found most commonly among gram-negative bacteria:
Escherichia coli (in some strains)
Neisseria meningitidis
Klebsiella pneumoniae
Haemophilus influenzae
Pseudomonas aeruginosa
Salmonella
Acinetobacter baumannii
However, some gram-positive bacteria may also have a capsule:
Bacillus megaterium for example, synthesizes a capsule composed of polypeptide and polysaccharides.
Bacillus anthracis
Streptococcus pyogenes synthesizes a hyaluronic acid capsule.
Streptococcus pneumoniae has at least 91 different capsular serotypes. These serotypes are the basis for the pneumococcal vaccines.
Streptococcus agalactiae produces a polysaccharide capsule of nine antigenic types that all contain sialic acid (Ia, Ib, II, III, IV, V, VI, VII, VIII).
Staphylococcus epidermidis
Staphylococcus aureus
Lactococcus garvieae synthesizes capsular gene clusters and some time synthesizes a hyaluronic acid capsule.
The yeast Cryptococcus neoformans, though not a bacterium, has a similar capsule.
Capsules too small to be seen with an ordinary microscope, such as the M protein of Streptococcus pyogenes, are called microcapsules.
Demonstration of capsule
India ink staining: the capsule appears as a clear halo around the bacterium as the ink can't penetrate the capsule.
Maneval's capsule stain: the capsule appears as a clear halo between the pink-stained bacterium and the bluish-grey stained background. The background stain is the acidic stain Congo red (which changes color to bluish-grey due to the pH), and the pink stain is fuchsine.
Serological methods: Capsular material is antigenic and can be demonstrated by mixing it with a specific anticapsular serum. When examined under the microscope, the capsule appears 'swollen' due to an increase in its refractivity. This phenomenon is the basis of quellung reaction.
Use in vaccination
Vaccination using capsular material is effective against some organisms (e.g., H. influenzae type b, S. pneumoniae, and N. meningitidis). However, polysaccharides are not highly antigenic, especially in children, so many capsular vaccines contain polysaccharides conjugated with protein carriers, such as the tetanus toxoid or diphtheria toxoid. This stimulates a much more robust immune response.
| Biology and health sciences | Basic anatomy | Biology |
1368943 | https://en.wikipedia.org/wiki/Mangla%20Dam | Mangla Dam | The Mangla Dam () is a multipurpose dam situated on the Jhelum River, lying in the Mirpur District of Azad Kashmir and the Jhelum District in Punjab, Pakistan. It is the sixth-largest dam in the world. The village of Mangla, which sits at the mouth of the dam, serves as its namesake. In November 1961, the project's selected contractors were revealed; it was announced that Binnie & Partners, a British engineering firm, was going to serve as the lead designers, engineers, and inspectors for the construction of the dam (led by Geoffrey Binnie). The project was undertaken by a consortium known as the Mangla Dam Contractors, which consisted of eight American construction firms sponsored by the Guy F. Atkinson Company based in South San Francisco, California.
Background
As part of the Indus Waters Treaty signed in 1960, India gained rights to the waters of the Ravi, Sutlej and Beas rivers, while Pakistan, in addition to the waters of the aforementioned three rivers' sections within Pakistani territory and some monetary compensation, received the rights to develop the Jhelum, Chenab and Indus river basins. Until 1967, the entire irrigation system of Pakistan was fully dependent on unregulated flows of the Indus River and its major tributaries. The agricultural yield was very low for a number of reasons, the most significant being a lack of water during critical growing periods. This problem stemmed from the seasonal variations in river flow due to monsoons and the absence of storage reservoirs to conserve the vast amounts of surplus water during those periods of high river discharge.
The Mangla Dam was the first of the two dams constructed to reduce this shortcoming and strengthen the irrigation system of the country as part of the Indus Basin Project, with the other being the Tarbela Dam situated on the Indus River in Swabi, Khyber Pakhtunkhwa.
Construction
Cost
The Mangla Dam was constructed at a cost of billion ( billion) with funding being provided by the World Bank and the Asian Development Bank.
Reservoir
The dam was constructed between 1961 and 1965 across the Jhelum River and Poonch River in the Mirpur District of Kashmir, about southeast of the capital city of Islamabad. The Mangla Dam components include a reservoir, main embankment, intake embankment, main spillway, emergency spillway, intake structures, 5 tunnels, and a power station. Besides the main dam, a dyke called Sukian – in length and a small dam called Jari Dam to block the Jari Nala – about beyond the new Mirpur town had to be constructed.
There was a total of 120 x 106 cubic yards (cu yds) of excavation for the reservoir whereas the total fill amounted to 142 x 106 cu yds and concrete to 1.96 x 106 cu yds respectively. The main embankment is earthfill with clay as the core material. Gravel and A-type sandstone are applied on the shoulders. The maximum height of embankment above the core trench is and the length is . The intake embankment is earthfill type with B-type sandstone as the core material. Gravel is applied on the shoulders. The maximum height of the intake embankment above the core trench is and the length is .
Sukian Dam is earthfill with B-type sandstone as the core material. A-type sandstone is applied on the shoulders. The maximum height of the intake embankment above the core trench is and the length is .
Jari Dam is also an earthfill type with silt as the core material. Gravel is applied on the shoulders of the dam. The maximum height of Jari dam above the core trench is and the length is . The main spillway is a submerged orifice type with 9 radial gates, x each; it has a maximum capacity of 1.1 million cusecs. The emergency spillway is a weir type with an erodible bund and a maximum capacity of 0.23 million cusecs. The 5 tunnels are steel and concrete lined and 1,560 feet long in bedrock. The internal diameter ranges between and .
Power house
The powerhouse, which consists of turbines, generators, and transformers, has been constructed at the toe of an intake embankment at the ground surface elevation of 865 feet SPD. The water to the powerhouse is supplied through five steel-lined tunnels of 30/26 feet diameter. Each tunnel is designed to feed two generating units. The powerhouse tailrace discharges into New Bong Canal, which has a length of 25,000 feet with a discharge capacity of about 49,000 cusecs, and terminates at an automatic gate control headworks at about 12 km downstream located near old Bong Escape Headworks.
There are ten vertical Francis type turbines in the powerhouse. Each of these turbines has an output of 138,00 bhp with a rated head of 295 feet of water. The first four turbines were manufactured by Mitsubishi Electric, Japan and were installed in 1969, turbines 5 and 6 are manufactured by ČKD Blansko, the Czech Republic and were installed in 1974, turbines 7-8 were manufactured by ACEC, Belgium and were installed in 1981, while the remaining two turbines are a make of Škoda, Czech Republic and were commissioned in 1994.
These turbines are connected to umbrella-type generators which have a generation capacity of 100 MW. Hitachi, Japan had provided generators for turbines 1–4 and 7-8 while Škoda generators are connected to turbines 5-6 and 9–10.
These generators are in turn connected to three-phase transformers. The transformers connected to turbines 1, 4, and 7 were manufactured by the Italy-based Savigliano. The transformers for turbine 5 and 6 are a make of Italtrafo, another Italian company, while the remaining five transformers were provided by Škoda.
Displacement & Resettlement
The Government of Pakistan had agreed to pay royalties to the Government of AJK (Azad Jammu and Kashmir) for the use of the water and electricity generated by the dam. Pakistan initially committed to supplying free electricity to the entire region of Azad Kashmir and providing complimentary clean water to the city of Mirpur. However, over time, this agreement was not transparently upheld, resulting in Pakistan encountering challenges in fulfilling its financial obligations related to royalty payments for electricity and water services in Azad Kashmir. Over 280 villages and the towns of Mirpur and Dadyal were submerged and over 110,000 people were displaced from the area as a result of the dam being built. Some of those affected by the dam were given work permits for Britain by the Government of Pakistan, and as a result, in many cities in the UK the majority of the Pakistani community originates from the Dadyal-Mirpur area of Kashmir. There are 747,000 Mirpuris in the United Kingdom, and the British Mirpuri community forms about 70% of the British Pakistani community. The percentage is greater in northern cities and towns. In Bradford, an industrial town in north-west England, it is estimated that roughly three-quarters of the population is from Mirpur, with a sizeable population also in Birmingham. At the time, many took up work in the textile and steel mills, due to the acute shortage of workers in England.
Operation
The project was designed primarily to increase the amount of water that could be used for irrigation from the flow of the Jhelum and its tributaries. Its secondary function was to generate electrical power from the irrigation releases at the artificial head of the reservoir. The project, though not initially designed as one, also works as a flood control structure by retaining water during the flood-prone season of Monsoon.
On 5 December 1971, the dam was damaged due to a bombing raid conducted by the Indian Air Force during the Indo-Pakistani War of 1971. This was against the international convention that large water reservoirs would not be targeted in war. As a consequence, the hydro project was temporarily out of service.
From the data available in 2009, the project had generated 183.551 billion units of low-cost Hydel energy since its commissioning. The annual generation during 2008-2009 was 4797.425 Million KWh while the station shared a peak load of 1150 MW which was 8.18% of the total WAPDA system peak.
On 1 September 2013, the water level in Mangla Dam reached a record height of 1237.15 feet against the maximum conservation level of 1242 feet. Radio Pakistan reported that "the water level in Mangla Dam has attained the maximum height of 1237.15 feet in the history and it is still increasing."
Mangla Dam Raising Project
Mangla reservoir had an initial reservoir capacity of , which reduced to in 2005 due to the sedimentation & was likely to reduce further. To counteract this phenomenon, the Mangla Dam Raising Project was started in 2004 and the main dam, spillway and its allied works were completed in 2009 at a cost of Rs. 101.384 billion. This project effectively raised the dam height by 30 feet to 482 feet (147 m), thereby raising the maximum water conservation level from 1202 feet to 1242 feet. This increased the dam's storage capacity from . Besides, it is expected that after raising the height of the Mangla Dam by 30 feet, the power house will generate 12 percent additional energy per year which will increase its installed capacity from 1,000 MW to 1,120 MW.
The Mangla Dam Raising Project, however, has affected more than 40,000 people living in the vicinity of the dam. The total cost of compensation and resettlement was Rs. 70 billion. The resettlement project includes the construction of New Mirpur City, four satellite towns (Islamgarh, Chakswari, Dadyal, Siakh) with all civic amenities, the Mirpur Bypass and two bridges across River Jehlum and Bong Canal respectively.
Mangla Power House Expansion Project
The dam was expanded in the era of Pervez Musharraf but it did not enhance the capacity of electric generation except increasing the level of water in the dam.
In November 2012, the United States announced a grant of $150 million for the expansion of the Mangla Dam powerhouse. Under the project, $400 million would be spent on the Mangla Dam powerhouse which is estimated to provide additional production for the next 40 years. The project, when complete, will increase the power generation capacity of the Mangla Dam to 1,310 MW from the existing 1,000 MW capacity. WAPDA successfully commissioned two refurbished generating units of the Mangla Hydel Power Station on 23 May 2022, increasing their capacity from 200 MW to 270 MW.
| Technology | Dams | null |
1369226 | https://en.wikipedia.org/wiki/Brix | Brix | Degrees Brix (symbol °Bx) is a measure of the dissolved solids in a liquid, and is commonly used to measure dissolved sugar content of a solution. One degree Brix is 1 gram of sucrose in 100 grams of solution and represents the strength of the solution as percentage by mass. If the solution contains dissolved solids other than pure sucrose, then the °Bx only approximates the dissolved solid content. For example, when one adds equal amounts of salt and sugar to equal amounts of water, the degrees of refraction (BRIX) of the salt solution rises faster than the sugar solution. The °Bx is traditionally used in the wine, sugar, carbonated beverage, fruit juice, fresh produce, maple syrup, and honey industries. The °Bx is also used for measuring the concentration of a cutting fluid mixed in water for metalworking processes.
Comparable scales for indicating sucrose content are: the Plato scale (°P), which is widely used by the brewing industry; the Oechsle scale used in German and Swiss wine making industries, amongst others; and the Balling scale, which is the oldest of the three systems and therefore mostly found in older textbooks, but is still in use in some parts of the world.
A sucrose solution with an apparent specific gravity (20°/20 °C) of 1.040 would be 9.99325 °Bx or 9.99359 °P while the representative sugar body, the International Commission for Uniform Methods of Sugar Analysis (ICUMSA), which favours the use of mass fraction, would report the solution strength as 9.99249%. Because the differences between the systems are of little practical significance (the differences are less than the precision of most common instruments) and wide historical use of the Brix unit, modern instruments calculate mass fraction using ICUMSA official formulas but report the result as °Bx.
Background
In the early 1800s, Karl Balling, followed by Adolf Brix, and finally the Normal-Commissions under Fritz Plato, prepared pure sucrose solutions of known strength, measured their specific gravities and prepared tables of percent sucrose by mass vs. measured specific gravity. Balling measured specific gravity to 3 decimal places, Brix to 5, and the Normal-Eichungs Kommission to 6 with the goal of the Commission being to correct errors in the 5th and 6th decimal place in the Brix table.
Equipped with one of these tables, a brewer wishing to know how much sugar was in his wort could measure its specific gravity and enter that specific gravity into the Plato table to obtain °Plato, which is the concentration of sucrose by percentage mass. Similarly, a vintner could enter the specific gravity of his must into the Brix table to obtain the °Bx, which is the concentration of sucrose by percent mass. It is important to point out that neither wort nor must is a solution of pure sucrose in pure water. Many other compounds are dissolved as well but these are either sugars, which behave similar to sucrose with respect to specific gravity as a function of concentration, or compounds that are present in small amounts (minerals, hop acids in wort, tannins, acids in must). In any case, even if °Bx is not representative of the exact amount of sugar in a must or fruit juice, it can be used for comparison of relative sugar content.
Measurement
Specific gravity
As specific gravity was the basis for the Balling, Brix and Plato tables, dissolved sugar content was originally estimated by measurement of specific gravity using a hydrometer or pycnometer. In modern times, hydrometers are still widely used, but where greater accuracy is required, an electronic oscillating U-tube meter may be employed. Whichever means is used, the analyst enters the tables with specific gravity and takes out (using interpolation if necessary) the sugar content in percent by mass.
If the analyst uses the Plato tables (maintained by the American Society of Brewing Chemists) they reports in °P. If using the Brix table (the current version of which is maintained by NIST and can be found on their website), they reports in °Bx. If using the ICUMSA tables, they would report in mass fraction (m.f.).
It is not, typically, actually necessary to consult tables as the tabulated °Bx or °P value can be printed directly on the hydrometer scale next to the tabulated value of specific gravity or stored in the memory of the electronic U-tube meter or calculated from polynomial fits to the tabulated data, in fact, the ICUMSA tables are calculated from a best-fit polynomial.
Also note that the tables in use today are not those published by Brix or Plato. Those workers measured true specific gravity reference to water at 4 °C using, respectively, 17.5 °C and 20 °C, as the temperature at which the density of a sucrose solution was measured. Both NBS and ASBC converted to apparent specific gravity at 20 °C/20 °C. The ICUMSA tables are based on more recent measurements on sucrose, fructose, glucose and invert sugar, and they tabulate true density and weight in air at 20 °C against mass fraction.
Refractive index
Dissolution of sucrose and other sugars in water changes not only its specific gravity but its optical properties, in particular its refractive index and the extent to which it rotates the plane of linearly polarized light. The refractive index, nD, for sucrose solutions of various percentage by mass has been measured and tables of nD vs. °Bx published. As with the hydrometer, it is possible to use these tables to calibrate a refractometer so that it reads directly in °Bx. Calibration is usually based on the ICUMSA tables, but the user of an electronic refractometer should verify this.
Infrared absorption
Sugars also have known infrared absorption spectra and this has made it possible to develop instruments for measuring sugar concentration using mid-infrared (MIR), non-dispersive infrared (NDIR), and Fourier transform infrared (FT-IR) techniques. In-line instruments are available that allow constant monitoring of sugar content in sugar refineries, beverage plants, wineries, etc. As with any other instruments, MIR and FT-IR instruments can be calibrated against pure sucrose solutions and thus report in °Bx, but there are other possibilities with these technologies, as they have the potential to distinguish between sugars and interfering substances. Newer MIR and NDIR instruments have up to five analyzing channels that allow corrections for interference between ingredients.
Tables
Specific gravity
Formulas derived from the NBS table above:
Approximate (R2=0.999 98) values can be computed from:
where SG is the apparent specific gravity of the solution at 20 °C/20 °C.
More accurate (R2=0.999 999 994) values are available from the polynomial:
,
RMS disagreement between the polynomial and the NBS table is 0.0009 °Bx.
Another accurate (R2=0.999 999 97) and simpler formula is:
The above formulas should not be used outside the range 1.00000 to 1.17874 SG (0 to 40 °Bx).
The Plato scale can be approximated with a mean average error of less than 0.02°P with the following equation:
or with even higher accuracy (average error less than 0.00053°P with respect to the ASBC tables) from the best-fit polynomial:
.
The difference between the °Bx and °P as calculated from the respective polynomials is:
The difference is generally less than ±0.0005 °Bx or °P with the exception being for weak solutions. As 0 °Bx is approached °P tend toward as much as 0.002 °P higher than the °Bx calculated for the same specific gravity.
Disagreements of this order of magnitude can be expected as the NBS and the ASBC used slightly different values for the density of air and pure water in their calculations for converting to apparent specific gravity. It should be clear from these comments that Plato and Brix are, for all but the most exacting applications, the same.
The ICUMSA polynomials are generally only published in the form where mass fraction is used to derive the density. As a result, they are omitted from this section.
Refractive index
When a refractometer is used, the Brix value can be obtained from the polynomial fit to the ICUMSA table:
,
where is the refractive index measured at the wavelength of the sodium D line (589.3 nm) at 20 °C. Temperature is important as refractive index changes dramatically with temperature. Many refractometers have built in "Automatic Temperature Compensation" (ATC), which is based on knowledge of the way the refractive index of sucrose changes. For example, the refractive index of a sucrose solution of strength less than 10 °Bx is such that a 1 °C change in temperature would cause the Brix reading to shift by about 0.06 °Bx. Beer, conversely, exhibits a change with temperature about three times this much. It is important, therefore, that users of refractometers either make sure the sample and prism of the instrument are both close to 20 °C or, if that is difficult to ensure, readings should be taken at 2 temperatures separated by a few degrees, the change per degree noted and the final recorded value referenced to 20 °C using the Bx vs. Temp slope information.
As solutes other than sucrose may affect the refractive index and the specific gravity differently, this refractive "Brix" value is not interchangeable with the traditional hydrometer Brix unless corrections are applied. The formal term for such a refractive value is "Refractometric Dry Substance" (RDS). See below.
Usage
The four scales are often used interchangeably since the differences are minor.
Brix is primarily used in fruit juice, wine making, carbonated beverage industry, starch and the sugar industry.
Plato is primarily used in brewing.
Balling appears on older saccharimeters and is still used in the South African wine industry and in some breweries.
Oechsle as direct reading for the sugar content is primary used in wine making in Germany, Switzerland and Luxembourg.
Brix is used in the food industry for measuring the approximate amount of sugars in fruits, vegetables, juices, wine, soft drinks and in the starch and sugar manufacturing industry. Different countries use the scales in different industries: In brewing, the UK uses specific gravity X 1000; Europe uses Plato degrees; and the US use a mix of specific gravity, degrees Brix, degrees Baumé, and degrees Plato. For fruit juices, 1.0 degree Brix is denoted as 1.0% sugar by mass. This usually correlates well with perceived sweetness.
Brix measurements are also used in the dairy industry to measure the quality of colostrum given to newborn calves, goats, and sheep.
Modern optical Brix meters are divided into two categories. In the first are the Abbe-based instruments in which a drop of the sample solution is placed on a prism; the result is observed through an eyepiece. The critical angle (the angle beyond which light is totally reflected back into the sample) is a function of the refractive index and the operator detects this critical angle by noting where a dark-bright boundary falls on an engraved scale. The scale can be calibrated in Brix or refractive index. Often the prism mount contains a thermometer that can be used to correct to 20 °C in situations where measurement cannot be made at exactly that temperature. These instruments are available in bench and handheld versions.
Digital refractometers also find the critical angle, but the light path is entirely internal to the prism. A drop of sample is placed on its surface, so the critical light beam never penetrates the sample. This makes it easier to read turbid samples. The light/dark boundary, whose position is proportional to the critical angle, is sensed by a CCD array. These meters are also available in bench top (laboratory) and portable (pocket) versions. This ability to easily measure Brix in the field makes it possible to determine ideal harvesting times of fruit and vegetables so that products arrive at the consumers in a perfect state or are ideal for subsequent processing steps such as vinification.
Due to higher accuracy and the ability to couple it with other measuring techniques (% and %alcohol), most soft drink companies and breweries use an oscillating U-tube density meter. Refractometers are still commonly used for fruit juice.
Brix and actual dissolved solids content
When a sugar solution is measured by refractometer or density meter, the °Bx or °P value obtained by entry into the appropriate table only represents the amount of dry solids dissolved in the sample if the dry solids are exclusively sucrose. This is seldom the case. Grape juice (must), for example, contains little sucrose but does contain glucose, fructose, acids, and other substances. In such cases, the °Bx value clearly cannot be equated with the sucrose content, but it may represent a good approximation to the total sugar content. For example, an 11.0% by mass D-Glucose ("grape sugar") solution measured 10.9 °Bx using a hand held instrument. For these reasons, the sugar content of a solution obtained by use of refractometry with the ICUMSA table is often reported as "Refractometric Dry Substance" (RDS), which could be thought of as an equivalent sucrose content. Where it is desirable to know the actual dry solids content, empirical correction formulas can be developed based on calibrations with solutions similar to those being tested. For example, in sugar refining, dissolved solids can be accurately estimated from refractive index measurement corrected by an optical rotation (polarization) measurement.
Alcohol has a higher refractive index (1.361) than water (1.333). As a consequence, a refractometer measurement made on a sugar solution once fermentation has begun results in a reading substantially higher than the actual solids content. Thus, an operator must be certain that the sample they are testing has not begun to ferment. (If fermentation has indeed started, a correction can be made by estimating alcohol concentration from the original, pre-fermentation reading, termed "OG" by homebrewers.) Brix or Plato measurements based on specific gravity are also affected by fermentation, but in the opposite direction; as ethanol is less dense than water, an ethanol/sugar/water solution gives a Brix or Plato reading that is artificially low.
| Physical sciences | Concentration | Basics and measurement |
20191692 | https://en.wikipedia.org/wiki/Helium%20atom | Helium atom | A helium atom is an atom of the chemical element helium. Helium is composed of two electrons bound by the electromagnetic force to a nucleus containing two protons along with two neutrons, depending on the isotope, held together by the strong force. Unlike for hydrogen, a closed-form solution to the Schrödinger equation for the helium atom has not been found. However, various approximations, such as the Hartree–Fock method, can be used to estimate the ground state energy and wavefunction of the atom. Historically, the first such helium spectrum calculation was done by Albrecht Unsöld in 1927. Its success was considered to be one of the earliest signs of validity of Schrödinger's wave mechanics.
Introduction
The quantum mechanical description of the helium atom is of special interest, because it is the simplest multi-electron system and can be used to understand the concept of quantum entanglement. The Hamiltonian of helium, considered as a three-body system of two electrons and a nucleus and after separating out the centre-of-mass motion, can be written as
where is the reduced mass of an electron with respect to the nucleus, and are the electron-nucleus distance vectors and . It is important to note that it operates not in normal space, but in a 6-dimensional configuration space . The nuclear charge, is 2 for helium. In the approximation of an infinitely heavy nucleus, we have and the mass polarization term disappears, so that in operator language, the Hamiltonian simplifies to:
The wavefunction belongs to the tensor product of combined spin states and combined spatial wavefunctions, and since this Hamiltonian only acts on spatial wavefunctions, we can neglect spin states until after solving the spatial wavefunction. This is possible since, for any general vector, one has that where is a combined spatial wavefunction and is the combined spin component. The Hamiltonian operator, since it only acts on the spatial component, gives the eigenvector equation:
which implies that one should find solutions for where is a general combined spatial wavefunction. This energy, however, is not degenerate with multiplicity given by the dimension of the space of combined spin states because of a symmetrization postulate, which requires that physical solutions for identical fermions should be totally antisymmetric, imposing a restriction on the choice of based on solutions . Hence the solutions are of the form: where is the energy eigenket spatial wavefunction and is a spin wavefunction such that is antisymmetric and is merely some superpostion of these states.
Since the Hamiltonian is independent of spin, it commutes with all spin operators. Since it is also rotationally invariant, the total x, y or z component of angular momentum operator also commutes with the Hamiltonian. From these commutation relations, and also commutes with the Hamiltonian which implies that energy is independent of and . Although the purely spatial form of the Hamiltonian implies that the energy is independent of , this would only be true in the absence of symmetrization postulate. Due to the symmetrization postulate, the choice of will influence the type of wavefunction required by symmetrization postulate which would in turn influence the energy of the state.
Other operators that commute with the Hamiltonian are the spatial exchange operator and the parity operator. However a good combination of mutually commuting operators are: , , , and . Hence the final solutions are given as:
where the energy is fold degenerate. For electrons, the total spin can have values of 0 or 1. A state with the quantum numbers: principal quantum number , total spin , angular quantum number and total angular momentum is denoted by .
States corresponding to , are called parahelium (singlet state, so called as there exist state) and are called orthohelium (triplet state, so called as there exist states). Since the spin exchange operator can be expressed in terms of dot product of spin vectors, eigenkets of spin exchange operators are also eigenkets of . Hence parahelium can also be said to be the spin anti-symmetric state (singlet state) or orthohelium to be spin symmetric state (triplet state).
The singlet state is given as:
and triplet states are given as:
as per symmetrization and total spin number requirement. It is observed that triplet states are symmetric and singlet states are antisymmetric. Since the total wavefunction is antisymmetric, a symmetric spatial wavefunction can only be paired with antisymmetric wavefunction and vice versa. Hence orthohelium (triplet state) has a symmetric spin wavefunction but an antisymmetric spatial wavefunction and parahelium (singlet state) has an antisymmetric spin wavefunction but a symmetric spatial wavefunction. Hence the type of wavefunction of each state is given above. The degeneracy solely comes from this spatial wavefunction. Note that for , there is no degeneracy in spatial wavefunction.
Alternatively, a more generalized representation of the above can be provided without considering the spatial and spin parts separately. This method is useful in situations where such manipulation is not possible, however, it can be applied wherever needed. Since the spin part is tensor product of spin Hilbert vector spaces, its basis can be represented by tensor product of each of the set, with each of the set, . Note that here but are in fact orthogonal. In the considered approximation (Pauli approximation), the wave function can be represented as a second order spinor with 4 components , where the indices describe the spin projection of both electrons in this coordinate system. The usual normalization condition, , follows from the orthogonality of all elements. This general spinor can be written as 2×2 matrix:
If the Hamiltonian had been spin dependent, we would not have been able to treat each these components independently as shown previously since the Hamiltonian need not act in the same manner for all four components.
The matrix can also be represented as a linear combination of any given basis of four orthogonal (in the vector-space of 2×2 matrices) constant matrices with scalar function coefficients as .
A convenient basis consists of one anti-symmetric matrix (with total spin , corresponding to a singlet state)and three symmetric matrices (with total spin , corresponding to a triplet state)
It is easy to show, that the singlet state is invariant under all rotations (a scalar entity), while the triplet are spherical vector tensor representations of an ordinary space vector , with the three components:
Since all spin interaction terms between the four components of in the above (scalar) Hamiltonian are neglected (e.g. an external magnetic field, or relativistic effects, like angular momentum coupling), the four Schrödinger equations can be solved independently. This is identical to the previously discussed method of finding spatial wavefunction eigenstates independently of the spin states, here spatial wavefunctions of different spin states correspond to the different components of the matrix.
The spin here only comes into play through the Pauli exclusion principle, which for fermions (like electrons) requires antisymmetry under simultaneous exchange of spin and coordinates (totally antisymmetric wavefunction condition)
Parahelium is then the singlet state with a symmetric spatial function and orthohelium is the triplet state with an antisymmetric spatial function .
Approximation methods
Following from the above approximation, effectively reducing three body problem to two body problem, we have:
This Hamiltonian for helium with two electrons can be written as a sum of two terms:
where the zero-order unperturbed Hamiltonian is
while the perturbation term:
is the electron-electron interaction. is just the sum of the two hydrogenic Hamiltonians:
whereare independent Coulomb field Hamiltonian of each electron. Since the unperturbed Hamiltonian is a sum of two independent Hamiltonians (i.e. are separable), the wavefunction must be of form where and are eigenkets of and respectively. However, the spatial wavefunction of the form need not correspond to physical states of identical electrons as per the symmetrization postulate. Thus, to obtain physical solutions symmetrization of the wavefunctions and is carried out.
The proper wave function then must be composed of the symmetric (+) and antisymmetric(−) linear combinations:or for the special cases of (both electrons have identical quantum numbers, parahelium only): .
This explains the absence of the state (with ) for orthohelium, where consequently (with ) is the metastable ground state.
Note that all wavefunction obtained thus far cannot be separated into wavefunctions of each particle (even for electrons with identical and where wavefunction is because then, the spin of the electrons are in a superposition of different spin states: and from ) i.e. the wavefunctions are always in superposition of some kind. In other words, one cannot completely determine states of particle 1 and 2, or measurements of all details, of each electrons cannot be made on one particle without affecting the other. This follows since the wavefunction is always a superposition of different states where each electron has unique . This is in agreement with Pauli exclusion principle.
We can infer from these wavefunctions that .
The corresponding energies are:
A good theoretical descriptions of helium including the perturbation term can be obtained within the Hartree–Fock and Thomas–Fermi approximations (see below).
The Hartree–Fock method is used for a variety of atomic systems. However it is just an approximation, and there are more accurate and efficient methods used today to solve atomic systems. The "many-body problem" for helium and other few electron systems can be solved quite accurately. For example, the ground state of helium is known to fifteen digits. In Hartree–Fock theory, the electrons are assumed to move in a potential created by the nucleus and the other electrons.
Ground state of Helium: Perturbation method
Since ground state corresponds to (1,0,0) state, there can only be one representation of such wavefunction whose spatial wavefunction is:
We note that the ground state energy of unperturbed helium atom as:Which is 30% larger than experimental data.
We can find the first order correction in energy due to electron repulsion in Hamiltonian :
The energy for ground state of helium in first order becomes compared to its experimental value of . A better approximation for ground state energy is obtained by choosing better trial wavefunction in variational method.
Screening effect
The energy that we obtained is too low because the repulsion term between the electrons was ignored, whose effect is to raise the energy levels. As gets bigger, our approach should yield better results, since the electron-electron repulsion term will get smaller.
is a central potential which is chosen so that the effect of the perturbation is small. The net effect of each electron on the motion of the other one is to screen somewhat the charge of the nucleus, so a simple guess for is
where is a screening constant and the quantity is the effective charge. The potential is a Coulomb interaction, so the corresponding individual electron energies are given by
and the corresponding spatial wave function is given by
If Ze was 1.70, that would make the expression above for the ground state energy agree with the experimental value E0 = −2.903 a.u. of the ground state energy of helium. Since in this case, the screening constant is S = 0.30. For the ground state of helium, for the average shielding approximation, the screening effect of each electron on the other one is equivalent to about of the electric charge.
Ground state of Helium: The variational method
To obtain a more accurate energy the variational principle can be applied to the electron-electron potential using the wave function
After integrating this, the result is:
This is closer to the experimental value, but if a better trial wave function is used, an even more accurate answer could be obtained. An ideal wave function would be one that doesn't ignore the influence of the other electron. In other words, each electron represents a cloud of negative charge which somewhat shields the nucleus so that the other electron actually sees an effective nuclear charge Z that is less than 2. A wave function of this type is given by:
Treating Z as a variational parameter to minimize H. The Hamiltonian using the wave function above is given by:
After calculating the expectation value of and Vee the expectation value of the Hamiltonian becomes:
The minimum value of Z needs to be calculated, so taking a derivative with respect to Z and setting the equation to 0 will give the minimum value of Z:
This shows that the other electron somewhat shields the nucleus reducing the effective charge from 2 to 1.69. This result matches with experimental results closely as well as the calculations of effective Z in screening effect. Hence, we obtain the most accurate result yet:
Where again, represents the ionization energy of hydrogen.
Perturbation theory for Helium
Consider the same setting where unperturbed Hamiltonian is:and perturbation is electron repulsion: .
In general, for (1s)(nl) state, in first order perturbation theory:with:where I is known as direct integral and J is known as exchange integral or exchange energy. If the combined spatial wavefunction is symmetric, its energy level has the + symbol in , whereas for the antisymmetric combined spatial wavefunction, has a minus symbol. Since due to the symmetrization postulate, the combined spatial wavefunctions differ on symmetric or antisymmetric nature, the J term is responsible for the splitting of energy levels between ortho and para helium states.
They are calculated as:
The first integral is said to be analogous to classical potential due to Coulomb interaction, where the squares of wavefunctions are interpreted as electron density. However, no such classical analog exists for the J term. Using Green's theorem, one can show that the J terms are always positive. From these, the diagram for energy level splitting can be roughly sketched. It also follows that for these states of helium, the energy of parallel spins cannot be more than that of antiparallel spins.
High precision theory
The Schrodinger equation for helium, like that of hydrogen, can be solved to accuracies equivalent to the most precise experimental values. Among the additional effects that must be included for these high accuracies include:
mass polarization: the dynamics of nucleus around the atomic center of mass.
relativity: Breit-Pauli corrections
quantum electrodynamics effects: the Lamb shift due to electron interaction with vacuum fluctuations.
Experimental value of ionization energy
Helium's first ionization energy is . This value was measured experimentally. The theoretic value of Helium atom's second ionization energy is . The total ground state energy of the helium atom is , or , which equals .
| Physical sciences | s-Block | Chemistry |
22726521 | https://en.wikipedia.org/wiki/Aquaculture%20of%20salmonids | Aquaculture of salmonids | The aquaculture of salmonids is the farming and harvesting of salmonid fish under controlled conditions for both commercial and recreational purposes. Salmonids (particularly salmon and rainbow trout), along with carp and tilapia, are the three most important fish groups in aquaculture. The most commonly commercially farmed salmonid is the Atlantic salmon (Salmo salar).
In the United States, Chinook salmon and rainbow trout are the most commonly farmed salmonids for recreational and subsistence fishing through the National Fish Hatchery System. In Europe, brown trout are the most commonly reared fish for recreational restocking. Commonly farmed non-salmonid fish groups include tilapia, catfish, black sea bass and bream. In 2007, the aquaculture of salmonids was worth USD $10.7 billion globally. Salmonid aquaculture production grew over ten-fold during the 25 years from 1982 to 2007. In 2012, the leading producers of salmonids were Norway, Chile, Scotland and Canada.
Much controversy exists about the ecological and health impacts of intensive salmonids aquaculture. Of particular concern are the impacts on wild salmon and other marine life.
Methods
The aquaculture or farming of salmonids can be contrasted with capturing wild salmonids using commercial fishing techniques. However, the concept of "wild" salmon as used by the Alaska Seafood Marketing Institute includes stock enhancement fish produced in hatcheries that have historically been considered ocean ranching. The percentage of the Alaska salmon harvest resulting from ocean ranching depends upon the species of salmon and location. Methods of salmonid aquaculture originated in late 18th-century fertilization trials in Europe. In the late 19th century, salmon hatcheries were used in Europe and North America. From the late 1950s, enhancement programs based on hatcheries were established in the United States, Canada, Japan, and the USSR. The contemporary technique using floating sea cages originated in Norway in the late 1960s.
Salmonids are usually farmed in two stages and in some places maybe more. First, the salmon are hatched from eggs and raised on land in freshwater tanks. Increasing the accumulated thermal units of water during incubation reduces time to hatching. When they are 12 to 18 months old, the smolt (juvenile salmon) are transferred to floating sea cages or net pens anchored in sheltered bays or fjords along a coast. This farming in a marine environment is known as mariculture. There they are fed pelleted feed for another 12 to 24 months, when they are harvested.
Norway produces 33% of the world's farmed salmonids, and Chile produces 31%. The coastlines of these countries have suitable water temperatures and many areas well protected from storms. Chile is close to large forage fisheries which supply fish meal for salmon aquaculture. Scotland and Canada are also significant producers; and it was reported in 2012 that the Norwegian government at that time controlled a significant fraction of the Canadian industry.
Modern salmonid farming systems are intensive. Their ownership is often under the control of huge agribusiness corporations, operating mechanized assembly lines on an industrial scale. In 2003, nearly half of the world’s farmed salmon was produced by just five companies.
Hatcheries
Modern commercial hatcheries for supplying salmon smolts to aquaculture net pens have been shifting to recirculating aquaculture systems (RAS)s where the water is recycled within the hatchery. This allows location of the hatchery to be independent of a significant fresh water supply and allows economical temperature control to both speed up and slow down the growth rate to match the needs of the net pens.
Conventional hatchery systems operate flow-through, where spring water or other water sources flow into the hatchery. The eggs are then hatched in trays and the salmon smolts are produced in raceways. The waste products from the growing salmon fry and the feed are usually discharged into the local river. Conventional flow-through hatcheries, for example the majority of Alaska's enhancement hatcheries, use more than of water to produce a kg of smolts.
An alternative method to hatching in freshwater tanks is to use spawning channels. These are artificial streams, usually parallel to an existing stream with concrete or rip-rap sides and gravel bottoms. Water from the adjacent stream is piped into the top of the channel, sometimes via a header pond to settle out sediment. Spawning success is often much better in channels than in adjacent streams due to the control of floods which in some years can wash out the natural redds. Because of the lack of floods, spawning channels must sometimes be cleaned out to remove accumulated sediment. The same floods which destroy natural redds also clean them out. Spawning channels preserve the natural selection of natural streams as no temptation exists, as in hatcheries, to use prophylactic chemicals to control diseases. However, exposing fish to wild parasites and pathogens using uncontrolled water supplies, combined with the high cost of spawning channels, makes this technology unsuitable for salmon aquaculture businesses. This type of technology is only useful for stock enhancement programs.
Sea cages
Sea cages, also called sea pens or net pens, are usually made of mesh framed with steel or plastic. They can be square or circular, across and deep, with volumes between . A large sea cage can contain up to 90,000 fish.
They are usually placed side by side to form a system called a seafarm or seasite, with a floating wharf and walkways along the net boundaries. Additional nets can also surround the seafarm to keep out predatory marine mammals. Stocking densities range from /m3 for Atlantic salmon and /m3 for Chinook salmon.
In contrast to closed or recirculating systems, the open net cages of salmonid farming lower production costs, but provide no effective barrier to the discharge of wastes, parasites, and disease into the surrounding coastal waters. Farmed salmon in open net cages can escape into wild habitats, for example, during storms.
An emerging wave in aquaculture is applying the same farming methods used for salmonids to other carnivorous finfish species, such as cod, bluefin tuna, halibut, and snapper. However, this is likely to have the same environmental drawbacks as salmon farming.
A second emerging wave in aquaculture is the development of copper alloys as netting materials. Copper alloys have become important netting materials because they are antimicrobial (i.e., they destroy bacteria, viruses, fungi, algae, and other microbes), so they prevent biofouling (i.e., the undesirable accumulation, adhesion, and growth of microorganisms, plants, algae, tubeworms, barnacles, mollusks, and other organisms). By inhibiting microbial growth, copper alloy aquaculture cages avoid costly net changes that are necessary with other materials. The resistance of organism growth on copper alloy nets also provides a cleaner and healthier environment for farmed fish to grow and thrive.
Feeding
With the amount of worldwide fish meal production being almost a constant amount for the last 30+ years and at maximum sustainable yield, much of the fish meal market has shifted from chicken and pig feed to fish and shrimp feeds as aquaculture has grown in this time.
Work continues on developing salmonid diet made from concentrated plant protein. As of 2014, an enzymatic process can be used to lower the carbohydrate content of barley, making it a high-protein fish feed suitable for salmon. Many other substitutions for fish meal are known, and diets containing zero fish meal are possible. For example, a planned closed-containment salmon fish farm in Scotland uses ragworms, algae, and amino acids as feed. Some of the eicosapentaenoic acid (EPA) and docosahexaenoic acid (DHA) in Omega-3 fatty acids may be replaced by land-based (non-marine) algae oil, reducing the harvest of wild fish as fish meal.
However, commercial economic animal diets are determined by least-cost linear programming models that are effectively competing with similar models for chicken and pig feeds for the same feed ingredients, and these models show that fish meal is more useful in aquatic diets than in chicken diets, where they can make the chickens taste like fish. Unfortunately, this substitution can result in lower levels of the highly valued omega-3 content in the farmed product. However, when vegetable oil is used in the growing diet as an energy source and a different finishing diet containing high omega-3 content fatty acids from either fish oil, algae oils, or some vegetable oils are used a few months before harvest, this problem is eliminated.
On a dry-dry basis, 2–4 kg of wild-caught fish are needed to produce 1 kg of salmon. The ratio may be reduced if non-fish sources are added. Wild salmon require about 10 kg of forage fish to produce 1 kg of salmon, as part of the normal trophic level energy transfer. The difference between the two numbers is related to farmed salmon feed containing other ingredients beyond fish meal and because farmed fish do not expend energy hunting.
In 2017 it was reported that the American company Cargill has been researching and developing alternative feeds with EWOS through its internal COMPASS programs in Norway, resulting in the proprietary RAPID feed blend. These methods studied macronutrient profiles of fish feed based upon geography and season. Using RAPID feed, salmon farms reduced the time to maturity of salmon to about 15 months, in a period one-fifth faster than usual.
Other feed additives
, 50-80% of the world fish oil production is fed to farmed salmonids.
Farm raised salmonids are also fed the carotenoids astaxanthin and canthaxanthin, so their flesh colour matches wild salmon, which also contain the same carotenoid pigments from their diet in the wild.
Harvesting
Modern harvesting methods are shifting towards using wet-well ships to transport live salmon to the processing plant. This allows the fish to be killed, bled, and filleted before rigor has occurred. This results in superior product quality to the customer, along with more humane processing. To obtain maximum quality, minimizing the level of stress is necessary in the live salmon until actually being electrically and percussively killed and the gills slit for bleeding. These improvements in processing time and freshness to the final customer are commercially significant and forcing the commercial wild fisheries to upgrade their processing to the benefit of all seafood consumers.
An older method of harvesting is to use a sweep net, which operates a bit like a purse seine net. The sweep net is a big net with weights along the bottom edge. It is stretched across the pen with the bottom edge extending to the bottom of the pen. Lines attached to the bottom corners are raised, herding some fish into the purse, where they are netted. Before killing, the fish are usually rendered unconscious in water saturated in carbon dioxide, although this practice is being phased out in some countries due to ethical and product quality concerns. More advanced systems use a percussive-stun harvest system that kills the fish instantly and humanely with a blow to the head from a pneumatic piston. They are then bled by cutting the gill arches and immediately immersing them in iced water. Harvesting and killing methods are designed to minimize scale loss, and avoid the fish releasing stress hormones, which negatively affect flesh quality.
Wild versus farmed
Wild salmonids are captured from wild habitats using commercial fishing techniques. Most wild salmonids are caught in North American, Japanese, and Russian fisheries. The following table shows the changes in production of wild salmonids and farmed salmonids over a period of 25 years, as reported by the FAO. Russia, Japan and Alaska all operate major hatchery based stock enhancement programs. The resulting fish hatchery fish are defined as "wild" for FAO and marketing purposes.
Issues
The US in their dietary guidelines for 2010 recommends eating 8 ounces per week of a variety of seafood and 12 ounces for lactating mothers, with no upper limits set and no restrictions on eating farmed or wild salmon.
In 2018, Canadian dietary guidelines recommended eating at least two servings of fish each week and choosing fish such as char, herring, mackerel, salmon, sardines, and trout.
Currently, much controversy exists about the ecological and health impacts of intensive salmonid aquaculture. Of particular concern are the impacts on wild salmonids and other marine life and on the incomes of commercial salmonid fishermen. However, the 'enhanced' production of salmon juveniles – which for instance lead to a double-digit proportion (20-50%) of the Alaska's yearly ‘wild’ salmon harvest - is not void of controversy, and the Alaska salmon harvest are highly dependent on the operation of Alaska’s Regional Aquaculture Associations. Furthermore, the sustainability of enhanced/hatchery-based ‘wild’ caught salmon has long been hotly debated, both from a scientific and political/marketing perspective. Such debate and positions were central to a 'halt' in the re-certification of Alaska salmon fisheries by the Marine Stewardship Council (MSC) in 2012. The Alaska salmon fisheries subsequently re-attained MSC-certification status; however the heavily hatchery-dependent Prince William Sound (PWS) unit of certification (“one of the most valuable fishing area in the State”) was for several years excluded from the MSC-certification (it remained ‘under assessment’ pending further analysis).
Disease and parasites
In 1972, Gyrodactylus, a monogenean parasite, was introduced with live trout and salmon from Sweden (Baltic stocks are resistant to it) into government-operated hatcheries in Norway. From the hatcheries, infected eggs, smolt, and fry were implanted in many rivers with the goal to strengthen the wild salmon stocks, but caused instead devastation to some of the wild salmon populations affected.
In 1984, infectious salmon anemia (ISAv) was discovered in Norway in an Atlantic salmon hatchery. Eighty percent of the fish in the outbreak died. ISAv, a viral disease, is now a major threat to the viability of Atlantic salmon farming. It is now the first of the diseases classified on List One of the European Commission’s fish health regimen. Amongst other measures, this requires the total eradication of the entire fish stock should an outbreak of the disease be confirmed on any farm. ISAv seriously affects salmon farms in Chile, Norway, Scotland, and Canada, causing major economic losses to infected farms. As the name implies, it causes severe anemia of infected fish. Unlike mammals, the red blood cells of fish have ribosomes, and can become infected with viruses. The fish develop pale gills, and may swim close to the water surface, gulping for air. However, the disease can also develop without the fish showing any external signs of illness, the fish maintain a normal appetite, and then they suddenly die. The disease can progress slowly throughout an infected farm, and in the worst cases, death rates may approach 100%. It is also a threat to the dwindling stocks of wild salmon. Management strategies include developing a vaccine and improving genetic resistance to the disease.
In the wild, diseases and parasites are normally at low levels, and kept in check by natural predation on weakened individuals. In crowded net pens, they can become epidemics. Diseases and parasites also transfer from farmed to wild salmon populations. A recent study in British Columbia links the spread of parasitic sea lice from river salmon farms to wild pink salmon in the same river. The European Commission (2002) concluded, "The reduction of wild salmonid abundance is also linked to other factors but there is more and more scientific evidence establishing a direct link between the number of lice-infested wild fish and the presence of cages in the same estuary." It is reported that wild salmon on the west coast of Canada are being driven to extinction by sea lice from nearby salmon farms. These predictions have been disputed by other scientists and recent harvests have indicated that the predictions were in error. In 2011, Scottish salmon farming introduced the use of farmed ballan wrasse for the purpose of cleaning farmed salmon of ectoparasites.
Globally, salmon production fell around 9% in 2015, in large part due to acute outbreaks of sea lice in Scotland and Norway. Lasers are used to reduce lice infections.
In the mid 1980s to the 1990s, bacterial kidney disease (BKD) caused by Renibacterium salmoninarum heavily impacted Chinook hatcheries in Idaho. The disease causes granulomatous inflammation that can lead to abscesses in the liver, spleen, and kidneys.
Pollution and contaminants
Salmonid farms are typically sited in marine ecosystems with good water quality, high water exchange rates, current speeds fast enough to prevent pollution of the bottom but slow enough to prevent pen damage, protection from major storms, reasonable water depth, and a reasonable distance from major infrastructure such as ports, processing plants, and logistical facilities such as airports. Logistical considerations are significant, and feed and maintenance labor must be transported to the facility and the product returned. Siting decisions are complicated by complex, politically driven permit problems in many countries that prevents optimal locations for the farms.
In sites without adequate currents, heavy metals can accumulate on the benthos (seafloor) near the salmon farms, particularly copper and zinc.
Contaminants are commonly found in the flesh of farmed and wild salmon. Health Canada in 2002 published measurements of PCBs, dioxins and furans and PDBEs in several varieties of fish. The farmed salmonids population had nearly 3 times the level of PCBs, more than 3 times the level of PDBEs, and nearly twice the level of dioxins and furans seen in the wild population. On the other hand, "Update of the monitoring of levels of dioxins and PCBs in food and feed", a 2012 study from the European Food Safety Authority, stated that farmed salmon and trout contained on average a many times lesser fraction of dioxins and PCBs than wild-caught salmon and trout."
A 2004 study, reported in Science, analysed farmed and wild salmon for organochlorine contaminants. They found the contaminants were higher in farmed salmon. Within the farmed salmon, European (particularly Scottish) salmon had the highest levels, and Chilean salmon the lowest. The FDA and Health Canada have established
a tolerance/limit for PCBs in commercial fish of 2000 ppb A follow-up study confirmed this, and found levels of dioxins, chlorinated pesticides, PCBs and other contaminants up to ten times greater in farmed salmon than wild Pacific salmon. On a positive note, further research using the same fish samples used in the previous study, showed that farmed salmon contained levels of beneficial fatty acids that were two to three times higher than wild salmon. A follow-up benefit-risk analysis on salmon consumption balanced the cancer risks with the (n–3) fatty acid advantages of salmon consumption. For this reason, current methods for this type of analysis take into consideration the lipid content of the sample in question. PCBs specifically are lipophilic, so are found in higher concentrations in fattier fish in general, thus the higher level of PCB in the farmed fish is in relation to the higher content of beneficial n–3 and n–6 lipids they contain. They found that recommended levels of (n-3) fatty acid consumption can be achieved eating farmed salmon with acceptable carcinogenic risks, but recommended levels of (n-3) EPA+DHA intake cannot be achieved solely from farmed (or wild) salmon without unacceptable carcinogenic risks. The conclusions of this paper from 2005 were that
"...consumers should not eat farmed fish from Scotland, Norway and eastern Canada more than three times a year; farmed fish from Maine, western Canada and Washington state no more than three to six times a year; and farmed fish from Chile no more than about six times a year. Wild chum salmon can be consumed safely as often as once a week, pink salmon, Sockeye and Coho about twice a month and Chinook just under once a month."
In 2005, Russia banned importing chilled fish from Norway, after samples of Norwegian farmed fish showed high levels of heavy metals. According to the Russian Minister of Agriculture Aleksey Gordeyev, levels of lead in the fish were 10 to 18 times higher than Russian safety standards and cadmium levels were almost four times higher.
Pollutants or toxins introduced by pisciculturists
In 2006, eight Norwegian salmon producers were caught in unauthorized and unlabeled use of nitrite in smoked and cured salmon. Norway applies EU regulations on food additives, according to which nitrite is allowed as a food additive in certain types of meat, but not fish. Fresh salmon was not affected.
Kurt Oddekalv, leader of the Green Warriors of Norway, argues that the scale of fish farming in Norway is unsustainable. Huge volumes of uneaten feed and fish excrement pollute the seabed, while chemicals designed to fight sea lice find their way into the food chain. He says: "If people knew this, they wouldn’t eat salmon", describing the farmed fish as "the most toxic food in the world". Don Stanifordthe former scientist turned activist/investigator and head of a small Global Alliance Against Industrial Aquacultureagrees, saying that a 10-fold increase in the use of some chemicals was seen in the 2016-2017 timeframe. The use of the toxic drug emamectin is rising fast. The levels of chemicals used to kill sea lice have breached environmental safety limits more than 100 times in the last 10 years.
Impact on wild salmonids
Farmed salmonids can, and often do, escape from sea cages. If the farmed salmonid is not native, it can compete with native wild species for food and habitat. If the farmed salmonid is native, it can interbreed with the wild native salmonids. Such interbreeding can reduce genetic diversity, disease resistance, and adaptability. In 2004, about 500,000 salmon and trout escaped from ocean net pens off Norway. Around Scotland, 600,000 salmon were released during storms. Commercial fishermen targeting wild salmon frequently catch escaped farm salmon. At one stage, in the Faroe Islands, 20 to 40 percent of all fish caught were escaped farm salmon. In 2017, about 263,000 farmed non-native Atlantic salmon escaped from a net in Washington waters in the 2017 Cypress Island Atlantic salmon pen break.
Sea lice, particularly Lepeophtheirus salmonis and various Caligus species, including C. clemensi and C. rogercresseyi, can cause deadly infestations of both farm-grown and wild salmon. Sea lice are naturally occurring and abundant ectoparasites which feed on mucus, blood, and skin, and migrate and latch onto the skin of salmon during planktonic nauplii and copepodid larval stages, which can persist for several days. Large numbers of highly populated, open-net salmon farms can create exceptionally large concentrations of sea lice; when exposed in river estuaries containing large numbers of open-net farms, many young wild salmon are infected, and do not survive as a result. Adult salmon may survive otherwise critical numbers of sea lice, but small, thin-skinned juvenile salmon migrating to sea are highly vulnerable. In 2007, mathematical studies of data available from the Pacific coast of Canada indicated the louse-induced mortality of pink salmon in some regions was over 80%. Later that year, in reaction to the 2007 mathematical study mentioned above, Canadian federal fisheries scientists Kenneth Brooks and Simon Jones published a critique titled "Perspectives on Pink Salmon and Sea Lice: Scientific Evidence Fails to Support the Extinction Hypothesis "
The time since these studies has shown a general increase in abundance of Pink Salmon in the Broughton Archipelago. Another comment in the scientific literature by Canadian Government Fisheries scientists Brian Riddell and Richard Beamish et al. came to the conclusion that there is no correlation between farmed salmon louse numbers and returns of pink salmon to the Broughton Archipelago. And in relation to the 2007 Krkosek extinction theory: "the data was [sic] used selectively and conclusions do not match with recent observations of returning salmon".
A 2008 meta-analysis of available data shows that salmonid farming reduces the survival of associated wild salmonid populations. This relationship has been shown to hold for Atlantic, steelhead, pink, chum, and coho salmon. The decrease in survival or abundance often exceeds 50%. However, these studies are all correlation analysis and correlation doesn't equal causation, especially when similar salmon declines were occurring in Oregon and California, which have no salmon aquaculture or marine net pens. Independent of the predictions of the failure of salmon runs in Canada indicated by these studies, the wild salmon run in 2010 was a record harvest.
A 2010 study that made the first use of sea lice count and fish production data from all salmon farms on the Broughton Archipelago found no correlation between the farm lice counts and wild salmon survival.
The authors conclude that the 2002 stock collapse was not caused by the farm sea lice population:
although the farm sea lice population during the out-migration of juvenile pink salmon was greater in 2000 than that of 2001,
there was a record salmon returning to spawn in 2001 (from the juveniles in 2000) compared with a 97% collapse in 2002 (from the juveniles in 2001).
The authors also note that initial studies had not investigated bacterial and viral causes for the event despite reports of bleeding at the base of the fins, a symptom often associated with infections, but not with sea lice exposure under laboratory conditions.
Wild salmon are anadromous. They spawn inland in fresh water and when young migrate to the ocean where they grow up. Most salmon return to the river where they were born, although some stray to other rivers. Concern exists about of the role of genetic diversity within salmon runs. The resilience of the population depends on some fish being able to survive environmental shocks, such as unusual temperature extremes. The effect of hatchery production on the genetic diversity of salmon is also unclear.
Genetic modification
Salmon have been genetically modified in laboratories so they can grow faster. A company, Aqua Bounty Farms, has developed a modified Atlantic salmon which grows nearly twice as fast (yielding a fully grown fish at 16–18 months rather than 30), and is more disease resistant, and cold tolerant. It also requires 10% less food. This was achieved using a chinook salmon gene sequence affecting growth hormones, and a promoter sequence from the ocean pout affecting antifreeze production. Normally, salmon produce growth hormones only in the presence of light. The modified salmon does not switch growth hormone production off. The company first submitted the salmon for FDA approval in 1996. In 2015, FDA has approved the AquAdvantage Salmon for commercial production. A concern with transgenic salmon is what might happen if they escape into the wild. One study, in a laboratory setting, found that modified salmon mixed with their wild cohorts were aggressive in competing, but ultimately failed.
Impact on wild predatory species
Sea cages can attract a variety of wild predators which can sometimes become entangled in associated netting, leading to injury or death. In Tasmania, Australian salmon-farming sea cages have entangled white-bellied sea eagles. This has prompted one company, Huon Aquaculture, to sponsor a bird rehabilitation centre and try more robust netting.
Ecological
Juvenile farmed Chinook have been shown to have higher rates of predation due to their larger size than wild juveniles upon release into marine environments. Their size correlates with the preferred size of prey for predators like birds, seals, and fish. This may have ecological implications because of the effect on feeding.
Impact on forage fish
The use of forage fish for fish meal production has been almost a constant for the last 30 years and at the maximum sustainable yield, while the market for fish meal has shifted from chicken, pig, and pet food to aquaculture diets. This market shift at constant production appears an economic decision implying that the development of salmon aquaculture had no impact on forage fish harvest rates.
Fish do not actually produce omega-3 fatty acids, but instead accumulate them from either consuming microalgae that produce these fatty acids, as is the case with forage fish like herring and sardines, or consuming forage fish, as is the case with fatty predatory fish like salmon. To satisfy this requirement, more than 50% of the world fish oil production is fed to farmed salmon.
In addition, salmon require nutritional intakes of protein, which is often supplied in the form of fish meal as the lowest-cost alternative. Consequently, farmed salmon consume more fish than they generate as a final product, though considerably more preferred as food.
Salmon Aquaculture Dialogue and ASC Salmon Standard
In 2004, the World Wide Fund for Nature (WWF)-USA initiated the Salmon Aquaculture Dialogue, one of several Aquaculture Dialogues. The aim of the dialogues was to produce an environmental and social standard for farmed salmon and other species (12 species currently, as of 2018). Since 2012, the standards elaborated by the multi-stakeholder Dialogues were passed-on to the Aquaculture Stewardship Council (ASC) which was created in 2010 to administer and developed them further. The first such standard was the ASC Salmon Standard (June 2012, and revised in 2017 after comprehensive public consultation). The WWF had originally identified what they called "seven key environmental and social impacts", characterised as:
Land-raised salmon
Recirculating aquaculture systems make it possible to farm salmon entirely on land, which as of 2019 is an ongoing initiative in the industry. However, large farmed salmon companies such as Mowi and Cermaq were not investing in such systems beyond the hatchery stage. In the United States, a major investor in the effort was Atlantic Sapphire, which plans to bring salmon raised in Florida to market in 2021. Other companies investing in the effort include Nordic Acquafarms and Whole Oceans.
Challenges
Nephrocalcinosis is an emerging issue in salmon aquaculture with the increased use of recirculating aquaculture systems. One key driver of nephrocalcinosis is exposure to high levels of CO2 in the rearing water.
Species
Atlantic salmon
In their natal streams, Atlantic salmon are considered a prized recreational fish, pursued by avid fly anglers during its annual runs. At one time, the species supported an important commercial fishery and a supplemental food fishery. However, the wild Atlantic salmon fishery is commercially dead; after extensive habitat damage and overfishing, wild fish make up only 0.5% of the Atlantic salmon available in world fish markets. The rest are farmed, predominantly from aquaculture in Chile, Canada, Norway, Russia, the United Kingdom, and Tasmania.
Atlantic salmon is, by far, the species most often chosen for farming. It is easy to handle, grows well in sea cages, commands a high market value, and adapts well to being farmed away from its native habitats.
Adult male and female fish are anesthetized. Eggs and sperm are "stripped", after the fish are cleaned and cloth dried. Sperm and eggs are mixed, washed, and placed into fresh water. Adults recover in flowing, clean, well-aerated water. Some researchers have studied cryopreservation of the eggs.
Fry are generally reared in large freshwater tanks for 12 to 20 months. Once the fish have reached the smolt phase, they are taken out to sea, where they are held for up to two years. During this time, the fish grow and mature in large cages off the coasts of Canada, the United States, or parts of Europe. Generally, cages are made of two nets; inner nets, which wrap around the cages, hold the salmon while outer nets, which are held by floats, keep predators out.
Many Atlantic salmon escape from cages at sea. Those salmon that further breed tend to lessen the genetic diversity of the species leading to lower survival rates, and lower catch rates. On the West Coast of North America, the non-native salmon could be an invasive threat, especially in Alaska and parts of Canada. This could cause them to compete with native salmon for resources. Extensive efforts are underway to prevent escapes and the potential spread of Atlantic salmon in the Pacific and elsewhere. The risk of Atlantic Salmon becoming a legitimate invasive threat on the Pacific Coast of N. America is questionable in light of both Canadian and American governments deliberately introducing this species by the millions for a 100-year period starting in the 1900s. Despite these deliberate attempts to establish this species on the Pacific coast; no established populations have been reported.
In 2007, 1,433,708 tonnes of Atlantic salmon were harvested worldwide with a value of $7.58 billion. Ten years later, in 2017, over 2 million tonnes of farmed Atlantic salmon were harvested.
Steelhead
In 1989, steelhead were reclassified into the Pacific trout as Oncorhynchus mykiss from the former binominals of Salmo gairdneri (Columbia River redband trout) and S. irideus (coastal rainbow trout). Steelhead are an anadromous form of rainbow trout that migrate between lakes and rivers and the ocean, and are also known as steelhead salmon or ocean trout.
Steelhead are raised in many countries throughout the world. Since the 1950s, production has grown exponentially, particularly in Europe and recently in Chile. Worldwide, in 2007, 604,695 tonnes of farmed steelhead were harvested, with a value of $2.59 billion. The largest producer is Chile. In Chile and Norway, the ocean-cage production of steelhead has expanded to supply export markets. Inland production of rainbow trout to supply domestic markets has increased strongly in countries such as Italy, France, Germany, Denmark, and Spain. Other significant producing countries include the United States, Iran, Germany, and the UK. Rainbow trout, including juvenile steelhead in fresh water, routinely feed on larval, pupal, and adult forms of aquatic insects (typically caddisflies, stoneflies, mayflies, and aquatic dipterana). They also eat fish eggs and adult forms of terrestrial insects (typically ants, beetles, grasshoppers, and crickets) that fall into the water. Other prey include small fish up to one-third of their length, crayfish, shrimp, and other crustaceans. As rainbow trout grow, the proportion of fish consumed increases in most populations. Some lake-dwelling forms may become planktonic feeders. In rivers and streams populated with other salmonid species, rainbow trout eat varied fish eggs, including those of salmon, brown and cutthroat trout, mountain whitefish, and the eggs of other rainbow trout. Rainbows also consume decomposing flesh from carcasses of other fish. Adult steelhead in the ocean feed primarily on other fish, squid, and amphipods. Cultured steelhead are fed a diet formulated to closely resemble their natural diet that includes fish meal, fish oil, vitamins and minerals, and the carotenoid asthaxanthin for pigmentation.
The steelhead is especially susceptible to enteric redmouth disease. Considerable research has been conducted on redmouth disease, as its implications for steelhead farmers are significant. The disease does not affect humans.
Coho salmon
The Coho salmon is the state animal of Chiba, Japan.
Coho salmon mature after only one year in the sea, so two separate broodstocks (spawners) are needed, alternating each year. Broodfish are selected from the salmon in the seasites and transferred to freshwater tanks for maturation and spawning.
Worldwide, in 2007, 115,376 tonnes of farmed Coho salmon were harvested with a value of $456 million. Chile, with about 90 percent of world production, is the primary producer with Japan and Canada producing the rest.
Chinook salmon
Chinook salmon are the state fish of Oregon, and are known as "king salmon" because of their large size and flavourful flesh. Those from the Copper River in Alaska are particularly known for their color, rich flavor, firm texture, and high omega-3 oil content. Alaska has a long-standing ban on finfish aquaculture that was enacted in 1989. (Alaska Stat. § 16.40.210)
Worldwide, in 2007, of farmed Chinook salmon were harvested with a value of $83 million. New Zealand is the largest producer of farmed king salmon, accounting for over half of world production (7,400 tonnes in 2005). Most of the salmon are farmed in the sea (mariculture) using a method sometimes called sea-cage ranching, which takes place in large floating net cages, about 25 m across and 15 m deep, moored to the sea floor in clean, fast-flowing coastal waters. Smolt (young fish) from freshwater hatcheries are transferred to cages containing several thousand salmon, and remain there for the rest of their lives. They are fed fishmeal pellets high in protein and oil.
Chinook salmon are also farmed in net cages placed in freshwater rivers or raceways, using techniques similar to those used for sea-farmed salmon. A unique form of freshwater salmon farming occurs in some hydroelectric canals in New Zealand. A site in Tekapo, fed by fast, cold waters from the Southern Alps, is the highest salmon farm in the world, above sea level.
Before they are killed, cage salmon are sometimes anaesthetised with a herbal extract. They are then spiked in the brain. The heart beats for a time as the animal is bled from its sliced gills. This method of relaxing the salmon when it is killed produces firm, long-keeping flesh. Lack of disease in wild populations and low stocking densities used in the cages means that New Zealand salmon farmers do not use antibiotics and chemicals that are often needed elsewhere.
Timeline
1527: The life history of the Atlantic salmon is described by Hector Boece of the University of Aberdeen, Scotland.
1763: Fertilization trials for Atlantic salmon take place in Germany. Later biologists refined these in Scotland and France.
1854: Salmon spawing beds and rearing ponds built along the bank of a river by the Dohulla Fishery, Ballyconneely, Ireland.
1864: Hatchery raised Atlantic salmon fry were released in the River Plenty, Tasmania in a failed attempt to establish a population in Australia
1892: Hatchery raised Atlantic salmon fry were released in the Umkomass river in South Africa in a failed attempt to establish a population in Africa.
Late 19th century: Salmon hatcheries are used in Europe, North America, and Japan to enhance wild populations.
1961: Hatchery raised Atlantic salmon fry were released in the rivers of the Falkland Islands in a failed attempt to establish a population in the South Atlantic.
Late 1960s: First salmon farms established in Norway and Scotland.
1970: Hatchery raised Atlantic salmon fry were released in the rivers of the Kerguelen Islands in a failed attempt to establish a population in the Indian Ocean.
Early 1970s: Salmon farms established in North America.
1973: Rainbow trout was first bred in Thailand on Doi Inthanon, the highest mountain in Thailand, as part of the Royal Project, an initiative of King Bhumibol Adulyadej (Rama IX).
1975: Gyrodactylus, a small monogenean parasite, spreads from Norwegian hatcheries to wild salmon, probably by means of fishing gear, and devastates some wild salmon populations.
Late 1970s: Salmon farms established in Chile and New Zealand.
1984: Infectious salmon anemia, a viral disease, is discovered in a Norwegian salmon hatchery. Eighty percent of the involved fish die.
1985: Salmon farms established in Australia.
1987: First reports of escaped Atlantic salmon being caught in wild Pacific salmon fisheries.
1988: A storm hits the Faroe Islands releasing millions of Atlantic salmon.
1989: Furunculosis, a bacterial disease, spreads through Norwegian salmon farms and wild salmon.
1996: World farmed salmon production exceeds wild salmon harvest.
2007: A swarm of Pelagia noctiluca jellyfish wipes out a 100,000 fish salmon farm in Northern Ireland.
2019: The first salmon fish farm in the Middle East is established in the United Arab Emirates.
2021: Open-net salmon farming is banned in Tierra del Fuego, Argentina.
In popular culture
Chapter 14 of Paul Torday's 2007 novel Salmon Fishing in the Yemen includes a description of a visit to the "McSalmon Aqua Farms" where salmon are raised caged in a sea loch in Scotland.
| Technology | Aquaculture | null |
21280496 | https://en.wikipedia.org/wiki/Visual%20perception | Visual perception | Visual perception is the ability to interpret the surrounding environment through photopic vision (daytime vision), color vision, scotopic vision (night vision), and mesopic vision (twilight vision), using light in the visible spectrum reflected by objects in the environment. This is different from visual acuity, which refers to how clearly a person sees (for example "20/20 vision"). A person can have problems with visual perceptual processing even if they have 20/20 vision.
The resulting perception is also known as vision, sight, or eyesight (adjectives visual, optical, and ocular, respectively). The various physiological components involved in vision are referred to collectively as the visual system, and are the focus of much research in linguistics, psychology, cognitive science, neuroscience, and molecular biology, collectively referred to as vision science.
Visual system
In humans and a number of other mammals, light enters the eye through the cornea and is focused by the lens onto the retina, a light-sensitive membrane at the back of the eye. The retina serves as a transducer for the conversion of light into neuronal signals. This transduction is achieved by specialized photoreceptive cells of the retina, also known as the rods and cones, which detect the photons of light and respond by producing neural impulses. These signals are transmitted by the optic nerve, from the retina upstream to central ganglia in the brain. The lateral geniculate nucleus, which transmits the information to the visual cortex. Signals from the retina also travel directly from the retina to the superior colliculus.
The lateral geniculate nucleus sends signals to the primary visual cortex, also called striate cortex. Extrastriate cortex, also called visual association cortex is a set of cortical structures, that receive information from striate cortex, as well as each other. Recent descriptions of visual association cortex describe a division into two functional pathways, a ventral and a dorsal pathway. This conjecture is known as the two streams hypothesis.
The human visual system is generally believed to be sensitive to visible light in the range of wavelengths between 370 and 730 nanometers of the electromagnetic spectrum. However, some research suggests that humans can perceive light in wavelengths down to 340 nanometers (UV-A), especially the young. Under optimal conditions these limits of human perception can extend to 310 nm (UV) to 1100 nm (NIR).
Study
The major problem in visual perception is that what people see is not simply a translation of retinal stimuli (i.e., the image on the retina), with the brain altering the basic information taken in. Thus people interested in perception have long struggled to explain what visual processing does to create what is actually seen.
Early studies
There were two major ancient Greek schools, providing a primitive explanation of how vision works.
The first was the "emission theory" of vision which maintained that vision occurs when rays emanate from the eyes and are intercepted by visual objects. If an object was seen directly it was by 'means of rays' coming out of the eyes and again falling on the object. A refracted image was, however, seen by 'means of rays' as well, which came out of the eyes, traversed through the air, and after refraction, fell on the visible object which was sighted as the result of the movement of the rays from the eye. This theory was championed by scholars who were followers of Euclid's Optics and Ptolemy's Optics.
The second school advocated the so-called 'intromission' approach which sees vision as coming from something entering the eyes representative of the object. With its main propagator Aristotle (De Sensu), and his followers, this theory seems to have some contact with modern theories of what vision really is, but it remained only a speculation lacking any experimental foundation. (In eighteenth-century England, Isaac Newton, John Locke, and others, carried the intromission theory of vision forward by insisting that vision involved a process in which rays—composed of actual corporeal matter—emanated from seen objects and entered the seer's mind/sensorium through the eye's aperture.)
Both schools of thought relied upon the principle that "like is only known by like", and thus upon the notion that the eye was composed of some "internal fire" that interacted with the "external fire" of visible light and made vision possible. Plato makes this assertion in his dialogue Timaeus (45b and 46b), as does Empedocles (as reported by Aristotle in his De Sensu, DK frag. B17).
Alhazen (965 – 1040) carried out many investigations and experiments on visual perception, extended the work of Ptolemy on binocular vision, and commented on the anatomical works of Galen. He was the first person to explain that vision occurs when light bounces on an object and then is directed to one's eyes.
Leonardo da Vinci (1452–1519) is believed to be the first to recognize the special optical qualities of the eye. He wrote "The function of the human eye ... was described by a large number of authors in a certain way. But I found it to be completely different." His main experimental finding was that there is only a distinct and clear vision at the line of sight—the optical line that ends at the fovea. Although he did not use these words literally he actually is the father of the modern distinction between foveal and peripheral vision.
Isaac Newton (1642–1726/27) was the first to discover through experimentation, by isolating individual colors of the spectrum of light passing through a prism, that the visually perceived color of objects appeared due to the character of light the objects reflected, and that these divided colors could not be changed into any other color, which was contrary to scientific expectation of the day.
Unconscious inference
Hermann von Helmholtz is often credited with the first modern study of visual perception. Helmholtz examined the human eye and concluded that it was incapable of producing a high-quality image. Insufficient information seemed to make vision impossible. He, therefore, concluded that vision could only be the result of some form of "unconscious inference", coining that term in 1867. He proposed the brain was making assumptions and conclusions from incomplete data, based on previous experiences.
Inference requires prior experience of the world.
Examples of well-known assumptions, based on visual experience, are:
light comes from above;
objects are normally not viewed from below;
faces are seen (and recognized) upright;
closer objects can block the view of more distant objects, but not vice versa; and
figures (i.e., foreground objects) tend to have convex borders.
The study of visual illusions (cases when the inference process goes wrong) has yielded much insight into what sort of assumptions the visual system makes.
Another type of unconscious inference hypothesis (based on probabilities) has recently been revived in so-called Bayesian studies of visual perception. Proponents of this approach consider that the visual system performs some form of Bayesian inference to derive a perception from sensory data. However, it is not clear how proponents of this view derive, in principle, the relevant probabilities required by the Bayesian equation. Models based on this idea have been used to describe various visual perceptual functions, such as the perception of motion, the perception of depth, and figure-ground perception. The "wholly empirical theory of perception" is a related and newer approach that rationalizes visual perception without explicitly invoking Bayesian formalisms.
Gestalt theory
Gestalt psychologists working primarily in the 1930s and 1940s raised many of the research questions that are studied by vision scientists today.
The Gestalt Laws of Organization have guided the study of how people perceive visual components as organized patterns or wholes, instead of many different parts. "Gestalt" is a German word that partially translates to "configuration or pattern" along with "whole or emergent structure". According to this theory, there are eight main factors that determine how the visual system automatically groups elements into patterns: Proximity, Similarity, Closure, Symmetry, Common Fate (i.e. common motion), Continuity as well as Good Gestalt (pattern that is regular, simple, and orderly) and Past Experience.
Analysis of eye movement
During the 1960s, technical development permitted the continuous registration of eye movement during reading, in picture viewing, and later, in visual problem solving, and when headset-cameras became available, also during driving.
The picture to the right shows what may happen during the first two seconds of visual inspection. While the background is out of focus, representing the peripheral vision, the first eye movement goes to the boots of the man (just because they are very near the starting fixation and have a reasonable contrast). Eye movements serve the function of attentional selection, i.e., to select a fraction of all visual inputs for deeper processing by the brain.
The following fixations jump from face to face. They might even permit comparisons between faces.
It may be concluded that the icon face is a very attractive search icon within the peripheral field of vision. The foveal vision adds detailed information to the peripheral first impression.
It can also be noted that there are different types of eye movements: fixational eye movements (microsaccades, ocular drift, and tremor), vergence movements, saccadic movements and pursuit movements. Fixations are comparably static points where the eye rests. However, the eye is never completely still, and gaze position will drift. These drifts are in turn corrected by microsaccades, very small fixational eye movements. Vergence movements involve the cooperation of both eyes to allow for an image to fall on the same area of both retinas. This results in a single focused image. Saccadic movements is the type of eye movement that makes jumps from one position to another position and is used to rapidly scan a particular scene/image. Lastly, pursuit movement is smooth eye movement and is used to follow objects in motion.
Face and object recognition
There is considerable evidence that face and object recognition are accomplished by distinct systems. For example, prosopagnosic patients show deficits in face, but not object processing, while object agnosic patients (most notably, patient C.K.) show deficits in object processing with spared face processing. Behaviorally, it has been shown that faces, but not objects, are subject to inversion effects, leading to the claim that faces are "special". Further, face and object processing recruit distinct neural systems. Notably, some have argued that the apparent specialization of the human brain for face processing does not reflect true domain specificity, but rather a more general process of expert-level discrimination within a given class of stimulus, though this latter claim is the subject of substantial debate. Using fMRI and electrophysiology Doris Tsao and colleagues described brain regions and a mechanism for face recognition in macaque monkeys.
The inferotemporal cortex has a key role in the task of recognition and differentiation of different objects. A study by MIT shows that subset regions of the IT cortex are in charge of different objects. By selectively shutting off neural activity of many small areas of the cortex, the animal gets alternately unable to distinguish between certain particular pairments of objects. This shows that the IT cortex is divided into regions that respond to different and particular visual features. In a similar way, certain particular patches and regions of the cortex are more involved in face recognition than other object recognition.
Some studies tend to show that rather than the uniform global image, some particular features and regions of interest of the objects are key elements when the brain needs to recognise an object in an image. In this way, the human vision is vulnerable to small particular changes to the image, such as disrupting the edges of the object, modifying texture or any small change in a crucial region of the image.
Studies of people whose sight has been restored after a long blindness reveal that they cannot necessarily recognize objects and faces (as opposed to color, motion, and simple geometric shapes). Some hypothesize that being blind during childhood prevents some part of the visual system necessary for these higher-level tasks from developing properly. The general belief that a critical period lasts until age 5 or 6 was challenged by a 2007 study that found that older patients could improve these abilities with years of exposure.
Cognitive and computational approaches
In the 1970s, David Marr developed a multi-level theory of vision, which analyzed the process of vision at different levels of abstraction. In order to focus on the understanding of specific problems in vision, he identified three levels of analysis: the computational, algorithmic and implementational levels. Many vision scientists, including Tomaso Poggio, have embraced these levels of analysis and employed them to further characterize vision from a computational perspective.
The computational level addresses, at a high level of abstraction, the problems that the visual system must overcome. The algorithmic level attempts to identify the strategy that may be used to solve these problems. Finally, the implementational level attempts to explain how solutions to these problems are realized in neural circuitry.
Marr suggested that it is possible to investigate vision at any of these levels independently. Marr described vision as proceeding from a two-dimensional visual array (on the retina) to a three-dimensional description of the world as output. His stages of vision include:
A 2D or primal sketch of the scene, based on feature extraction of fundamental components of the scene, including edges, regions, etc. Note the similarity in concept to a pencil sketch drawn quickly by an artist as an impression.
A 2 D sketch of the scene, where textures are acknowledged, etc. Note the similarity in concept to the stage in drawing where an artist highlights or shades areas of a scene, to provide depth.
A 3 D model, where the scene is visualized in a continuous, 3-dimensional map.
Marr's 2D sketch assumes that a depth map is constructed, and that this map is the basis of 3D shape perception. However, both stereoscopic and pictorial perception, as well as monocular viewing, make clear that the perception of 3D shape precedes, and does not rely on, the perception of the depth of points. It is not clear how a preliminary depth map could, in principle, be constructed, nor how this would address the question of figure-ground organization, or grouping. The role of perceptual organizing constraints, overlooked by Marr, in the production of 3D shape percepts from binocularly-viewed 3D objects has been demonstrated empirically for the case of 3D wire objects, e.g. For a more detailed discussion, see Pizlo (2008).
A more recent, alternative framework proposes that vision is composed instead of the following three stages: encoding, selection, and decoding. Encoding is to sample and represent visual inputs (e.g., to represent visual inputs as neural activities in the retina). Selection, or attentional selection, is to select a tiny fraction of input information for further processing, e.g., by shifting gaze to an object or visual location to better process the visual signals at that location. Decoding is to infer or recognize the selected input signals, e.g., to recognize the object at the center of gaze as somebody's face. In this framework, attentional selection starts at the primary visual cortex along the visual pathway, and the attentional constraints impose a dichotomy between the central and peripheral visual fields for visual recognition or decoding.
Transduction
Transduction is the process through which energy from environmental stimuli is converted to neural activity. The retina contains three different cell layers: photoreceptor layer, bipolar cell layer, and ganglion cell layer. The photoreceptor layer where transduction occurs is farthest from the lens. It contains photoreceptors with different sensitivities called rods and cones. The cones are responsible for color perception and are of three distinct types labeled red, green, and blue. Rods are responsible for the perception of objects in low light. Photoreceptors contain within them a special chemical called a photopigment, which is embedded in the membrane of the lamellae; a single human rod contains approximately 10 million of them. The photopigment molecules consist of two parts: an opsin (a protein) and retinal (a lipid). There are 3 specific photopigments (each with their own wavelength sensitivity) that respond across the spectrum of visible light. When the appropriate wavelengths (those that the specific photopigment is sensitive to) hit the photoreceptor, the photopigment splits into two, which sends a signal to the bipolar cell layer, which in turn sends a signal to the ganglion cells, the axons of which form the optic nerve and transmit the information to the brain. If a particular cone type is missing or abnormal, due to a genetic anomaly, a color vision deficiency, sometimes called color blindness will occur.
Opponent process
Transduction involves chemical messages sent from the photoreceptors to the bipolar cells to the ganglion cells. Several photoreceptors may send their information to one ganglion cell. There are two types of ganglion cells: red/green and yellow/blue. These neurons constantly fire—even when not stimulated. The brain interprets different colors (and with a lot of information, an image) when the rate of firing of these neurons alters. Red light stimulates the red cone, which in turn stimulates the red/green ganglion cell. Likewise, green light stimulates the green cone, which stimulates the green/red ganglion cell and blue light stimulates the blue cone which stimulates the blue/yellow ganglion cell. The rate of firing of the ganglion cells is increased when it is signaled by one cone and decreased (inhibited) when it is signaled by the other cone. The first color in the name of the ganglion cell is the color that excites it and the second is the color that inhibits it. i.e.: A red cone would excite the red/green ganglion cell and the green cone would inhibit the red/green ganglion cell. This is an opponent process. If the rate of firing of a red/green ganglion cell is increased, the brain would know that the light was red, if the rate was decreased, the brain would know that the color of the light was green.
Artificial visual perception
Theories and observations of visual perception have been the main source of inspiration for computer vision (also called machine vision, or computational vision). Special hardware structures and software algorithms provide machines with the capability to interpret the images coming from a camera or a sensor.
| Biology and health sciences | Nervous system | null |
21280576 | https://en.wikipedia.org/wiki/Mechanism%20%28engineering%29 | Mechanism (engineering) | In engineering, a mechanism is a device that transforms input forces and movement into a desired set of output forces and movement. Mechanisms generally consist of moving components which may include Gears and gear trains; Belts and chain drives; cams and followers; Linkages; Friction devices, such as brakes or clutches; Structural components such as a frame, fasteners, bearings, springs, or lubricants; Various machine elements, such as splines, pins, or keys.
German scientist Franz Reuleaux defines machine as "a combination of resistant bodies so arranged that by their means the mechanical forces of nature can be compelled to do work accompanied by certain determinate motion". In this context, his use of machine is generally interpreted to mean mechanism.
The combination of force and movement defines power, and a mechanism manages power to achieve a desired set of forces and movement.
A mechanism is usually a piece of a larger process, known as a mechanical system or machine. Sometimes an entire machine may be referred to as a mechanism; examples are the steering mechanism in a car, or the winding mechanism of a wristwatch.
However, typically, a set of multiple mechanisms is called a machine.
Kinematic pairs
From the time of Archimedes to the Renaissance, mechanisms were viewed as constructed from simple machines, such as the lever, pulley, screw, wheel and axle, wedge, and inclined plane. Reuleaux focused on bodies, called links, and the connections between these bodies, called kinematic pairs, or joints.
To use geometry to study the movement of a mechanism, its links are modelled as rigid bodies. This means that distances between points in a link are assumed to not change as the mechanism moves—that is, the link does not flex. Thus, the relative movement between points in two connected links is considered to result from the kinematic pair that joins them.
Kinematic pairs, or joints, are considered to provide ideal constraints between two links, such as the constraint of a single point for pure rotation, or the constraint of a line for pure sliding, as well as pure rolling without slipping and point contact with slipping. A mechanism is modelled as an assembly of rigid links and kinematic pairs.
Links and joints
Reuleaux called the ideal connections between links kinematic pairs. He distinguished between higher pairs, with line contact between the two links, and lower pairs, with area contact between the links. shows that there are many ways to construct pairs that do not fit this simple model.
Lower pair: A lower pair is an ideal joint that has surface contact between the pair of elements, as in the following cases:
A revolute pair, or hinged joint, requires that a line in the moving body remain co-linear with a line in the fixed body, and a plane perpendicular to this line in the moving body must maintain contact with a similar perpendicular plane in the fixed body. This imposes five constraints on the relative movement of the links, which therefore gives the pair one degree of freedom.
A prismatic joint, or slider, requires that a line in the moving body remain co-linear with a line in the fixed body, and a plane parallel to this line in the moving body must maintain contact with a similar parallel plane in the fixed body. This imposes five constraints on the relative movement of the links, which therefore gives the pair one degree of freedom.
A cylindrical joint requires that a line in the moving body remain co-linear with a line in the fixed body. It combines a revolute joint and a sliding joint. This joint has two degrees of freedom.
A spherical joint, or ball joint, requires that a point in the moving body maintain contact with a point in the fixed body. This joint has three degrees of freedom.
A planar joint requires that a plane in the moving body maintain contact with a plane in a fixed body. This joint has three degrees of freedom.
A screw joint, or helical joint, has only one degree of freedom because the sliding and rotational motions are related by the helix angle of the thread.
Higher pairs: Generally, a higher pair is a constraint that requires a line or point contact between the elemental surfaces. For example, the contact between a cam and its follower is a higher pair called a cam joint. Similarly, the contact between the involute curves that form the meshing teeth of two gears are cam joints.
Kinematic diagram
A kinematic diagram reduces machine components to a skeleton diagram that emphasises the joints and reduces the links to simple geometric elements. This diagram can also be formulated as a graph by representing the links of the mechanism as edges and the joints as vertices of the graph. This version of the kinematic diagram has proven effective in enumerating kinematic structures in the process of machine design.
An important consideration in this design process is the degree of freedom of the system of links and joints, which is determined using the Chebychev–Grübler–Kutzbach criterion.
Planar mechanisms
While all mechanisms in a mechanical system are three-dimensional, they can be analysed using plane geometry if the movement of the individual components is constrained so that all point trajectories are parallel or in a series connection to a plane. In this case the system is called a planar mechanism. The kinematic analysis of planar mechanisms uses the subset of Special Euclidean group SE, consisting of planar rotations and translations, denoted by SE.
The group SE is three-dimensional, which means that every position of a body in the plane is defined by three parameters. The parameters are often the x and y coordinates of the origin of a coordinate frame in M, measured from the origin of a coordinate frame in F, and the angle measured from the x-axis in F to the x-axis in M. This is often described saying a body in the plane has three degrees of freedom.
The pure rotation of a hinge and the linear translation of a slider can be identified with subgroups of SE, and define the two joints as one degree-of-freedom joints of planar mechanisms. The cam joint formed by two surfaces in sliding and rotating contact is a two degree-of-freedom joint.
Spherical mechanisms
It is possible to construct a mechanism such that the point trajectories in all components lie in concentric spherical shells around a fixed point. An example is the gimbaled gyroscope. These devices are called spherical mechanisms. Spherical mechanisms are constructed by connecting links with hinged joints such that the axes of each hinge pass through the same point. This point becomes centre of the concentric spherical shells. The movement of these mechanisms is characterised by the group SO(3) of rotations in three-dimensional space. Other examples of spherical mechanisms are the automotive differential and the robotic wrist.
The rotation group SO(3) is three-dimensional. An example of the three parameters that specify a spatial rotation are the roll, pitch and yaw angles used to define the orientation of an aircraft.
Spatial mechanisms
A mechanism in which a body moves through a general spatial movement is called a spatial mechanism. An example is the RSSR linkage, which can be viewed as a four-bar linkage in which the hinged joints of the coupler link are replaced by rod ends, also called spherical joints or ball joints. The rod ends let the input and output cranks of the RSSR linkage be misaligned to the point that they lie in different planes, which causes the coupler link to move in a general spatial movement. Robot arms, Stewart platforms, and humanoid robotic systems are also examples of spatial mechanisms.
Bennett's linkage is an example of a spatial overconstrained mechanism, which is constructed from four hinged joints.
The group SE(3) is six-dimensional, which means the position of a body in space is defined by six parameters. Three of the parameters define the origin of the moving reference frame relative to the fixed frame. Three other parameters define the orientation of the moving frame relative to the fixed frame.
Linkages
A linkage is a collection of links connected by joints. Generally, the links are the structural elements and the joints allow movement. Perhaps the single most useful example is the planar four-bar linkage. There are, however, many more special linkages:
Watt's linkage is a four-bar linkage that generates an approximate straight line. It was critical to the operation of his design for the steam engine. This linkage also appears in vehicle suspensions to prevent side-to-side movement of the body relative to the wheels.
The success of Watt's linkage led to the design of similar approximate straight-line linkages, such as Hoeken's linkage and Chebyshev's linkage.
The Peaucellier linkage generates a true straight-line output from a rotary input.
The Sarrus linkage is a spatial linkage that generates straight-line movement from a rotary input.
The Klann linkage and the Jansen linkage are recent inventions that provide interesting walking movements. They are respectively a six-bar and an eight-bar linkage.
Compliant mechanisms
A compliant mechanism is a series of rigid bodies connected by compliant elements. These mechanisms have many advantages, including reduced part-count, reduced "slop" between joints (no parasitic motion because of gaps between parts), energy storage, low maintenance (they don't require lubrication and there is low mechanical wear), and ease of manufacture.
Flexure bearings (also known as flexure joints) are a subset of compliant mechanisms that produce a geometrically well-defined motion (rotation) on application of a force.
Cam and follower mechanisms
A cam and follower mechanism is formed by the direct contact of two specially shaped links. The driving link is called the cam and the link that is driven through the direct contact of their surfaces is called the follower. The shape of the contacting surfaces of the cam and follower determines the movement of the mechanism. In general a cam and follower mechanism's energy is transferred from cam to follower. The camshaft is rotated and, according to the cam profile, the follower moves up and down. Nowadays, slightly different types of eccentric cam followers are also available, in which energy is transferred from the follower to the cam. The main benefit of this type of cam and follower mechanism is that the follower moves slightly and helps to rotate the cam six times more circumference length with 70% of the force.
Gears and gear trains
The transmission of rotation between contacting toothed wheels can be traced back to the Antikythera mechanism of Greece and the south-pointing chariot of China. Illustrations by the Renaissance scientist Georgius Agricola show gear trains with cylindrical teeth. The implementation of the involute tooth yielded a standard gear design that provides a constant speed ratio. Some important features of gears and gear trains are:
The ratio of the pitch circles of mating gears defines the speed ratio and the mechanical advantage of the gear set.
A planetary gear train provides high gear reduction in a compact package.
It is possible to design gear teeth for gears that are non-circular, yet still transmit torque smoothly.
The speed ratios of chain and belt drives are computed in the same way as gear ratios.
Mechanism synthesis
The design of mechanisms to achieve a particular movement and force transmission is known as the kinematic synthesis of mechanisms. This is a set of geometric techniques which yield the dimensions of linkages, cam and follower mechanisms, and gears and gear trains to perform a required mechanical movement and power transmission.
| Technology | Machinery and tools: General | null |
28802456 | https://en.wikipedia.org/wiki/Lunar%20rover | Lunar rover | A lunar rover or Moon rover is a space exploration vehicle designed to move across the surface of the Moon. The Apollo program's Lunar Roving Vehicle was driven on the Moon by members of three American crews, Apollo 15, 16, and 17. Other rovers have been partially or fully autonomous robots, such as the Soviet Union's Lunokhods, Chinese Yutus, Indian Pragyan, and Japan's LEVs. Five countries have had operating rovers on the Moon: the Soviet Union, the United States, China, India, and Japan.
Variations in design
Lunar rover designs have varied in several ways.
Size and speed
Lunokhod rovers were in length. The LRVs were long with a wheelbase, and achieved a top speed of during Apollo 17.
Power
The Lunokhod rovers, and others, used photovoltaic solar power. The LRV rovers were battery powered. Lunokhod and the Chinese Yutu rovers were furthermore equipped with a radioisotope heater unit to keep instruments warm. These, however, delivered only heat, not electric power. While unlike on other celestial bodies, such as Earth or Mars, there is no atmosphere to interfere with solar power, the extreme length of the day/night cycle complicates the use of solar power as energy storage or hibernation are necessary for any missions exceeding two weeks in length. There are places where solar power is almost always available (especially near the lunar south pole) on the Moon, but to date no mission has successfully landed a rover at one of those places. Radioisotope thermoelectric generators can operate independent of the day/night cycle and have been used on missions to other celestial bodies in the past.
Propulsion
The LRV was a four-wheel design. The Lunokhod rovers used eight.
Thermal control
To remain warm during periods of lunar night the Lunokhod rovers used heat from radioactive polonium-210.
Past missions
Lunokhod 1
In November 1970, as part of the Lunokhod program, the Soviet Union sent the Lunokhod 1 robotic rover to the lunar surface. It remained operational until October 1971. The rover was soft-landed in Mare Imbrium by the Luna 17 lander. Lunokhod 1 was the first rover to land on another celestial body.
Apollo Lunar Roving Vehicle
The Lunar Roving Vehicle (LRV) was a battery-powered four-wheeled vehicle design. The LRV could carry one or two astronauts, their equipment, and lunar samples. During 1971 and 1972, LRVs were used on the Moon for each of the final three missions of the American Apollo program, Apollo 15, 16, and 17.
Lunokhod 2
Lunokhod 2 was the second of two monocrystalline-panel-powered uncrewed lunar rovers landed on the Moon by the Soviet Union as part of the Lunokhod program. The Luna 21 spacecraft landed on the Moon and deployed the second Soviet lunar rover Lunokhod 2 in January 1973. The objectives of the mission were to collect images of the lunar surface, examine ambient light levels to determine the feasibility of astronomical observations from the Moon, perform laser ranging experiments, observe solar X-rays, measure local magnetic fields, and study the soil mechanics of the lunar surface material. Lunokhod 2 was intended to be followed by Lunokhod 3 (No.205) in 1977 but the mission was cancelled.
Yutu
Yutu is a Chinese lunar rover that launched on 1 December 2013 and landed on 14 December 2013 as part of the Chang'e 3 mission. It is China's first lunar rover, part of the second phase of the Chinese Lunar Exploration Program undertaken by China National Space Administration (CNSA). The lunar rover is called Yutu, or Jade Rabbit, a name selected in an online poll.
The rover encountered operational difficulties after the first 14-day lunar night, and was unable to move after the end of the second lunar night, finally on August 3, 2016, it officially stopped sending data and doing its operations.
Pragyan (Chandrayaan-3 rover)
Chandrayaan-3 was launched on 14 July 2023 by the Indian Space Research Organisation in India's second attempt to soft land a rover and a lander on the Moon. Pragyan became the first rover to operate near the Moon's south pole when it successfully landed on 23 August 2023, after the lander separation from propulsion module had taken place on 17 August. The Pragyan rover was deployed the same day as landing and has travelled since then. On September 2, the rover finished all assignments and entered into a sleep mode in preparation for wake up on September 22, but was unable to do so.
SLIM's LEV Rovers
The SLIM lander has two rovers onboard, Lunar Excursion Vehicle 1 (LEV-1) (hopper) and Lunar Excursion Vehicle 2 (LEV-2), also known as Sora-Q, a tiny rover developed by JAXA in joint cooperation with Tomy, Sony Group, and Doshisha University. The first rover has direct-to-Earth communication. The second rover is designed to change its shape to traverse around the landing site over a short lifespan of two hours. SLIM was launched on September 6, 2023, and reached lunar orbit on 25 December 2023. The two rovers were successfully deployed and landed separately from SLIM shortly before its own landing on 19 January 2024. LEV-1 conducted seven hops over 107 minutes on lunar surface and LEV-2 imaged SLIM on lunar surface.
Jinchan
Chinese Chang'e 6 sample return mission carries a mini rover called Jinchan to conduct infrared spectroscopy of lunar surface and imaged Chang'e 6 lander on lunar surface.
Failed missions
Pragyan (Chandrayaan-2 rover)
Chandrayaan-2 was the second lunar mission by India, consisting of a lunar orbiter, a lander named Vikram, and a rover named Pragyan. The rover weighing 27 kg, had six wheels and was to be operated on solar power. Launched on 22 July 2019, the mission entered lunar orbit on August 20. Pragyan was destroyed along with its lander, Vikram, when it crash-landed on the Moon on 6 September 2019 and never got the chance to deploy.
Rashid
Rashid was a lunar rover built by MBRSC to be launched onboard Ispace's lander called Hakuto-R. The rover was launched in November 2022, but was destroyed as the lander crash landed in April 2023. It was equipped with two high-resolution cameras, a microscopic camera to capture small details, and a thermal imaging camera. The rover carried a Langmuir probe, designed to study the Moon's plasma and will attempt to explain why Moon dust is so sticky. The rover was supposed to study the lunar surface, mobility on the Moon’s surface and how different surfaces interact with lunar particles.
SORA-Q
Sora-Q was developed by Takara Tomy, JAXA and Doshisha University to be launched onboard Ispace's lander called Hakuto-R Mission 1. It was launched in 2022, but was destroyed as the lander crash landed in April 2023. A second rover was successfully deployed from the SLIM lander in January 2024.
Peregrine Mission One
Peregrine lander launched on 8 January 2024 to the Moon. It took with it 5 Colmena rovers and a Iris rover. The mission of the Peregrine lander was forced to be cancelled after an excessive propellant leak.
Active missions
Yutu-2
The Chang'e 4 Chinese mission launched on 7 December 2018, and landed and deployed the Yutu-2 rover on the far side of the Moon on 3 January 2019. It is the first rover to operate on the Moon's far side.
In December 2019, Yutu 2 broke the lunar longevity record, previously held by the Soviet Union's Lunokhod 1 rover, which operated on the lunar surface for eleven lunar days (321 Earth days) and traversed a total distance of .
In February 2020, Chinese astronomers reported, for the first time, a high-resolution image of a lunar ejecta sequence, and, as well, direct analysis of its internal architecture. These were based on observations made by the Lunar Penetrating Radar (LPR) on board the Yutu-2 rover while studying the far side of the Moon.
Data from its two-channel ground penetrating radar (GPR) has been used by scientists to put together an image of multiple layers beneath the surface of the far side of the Moon up to a depth of 300 meters.
Yutu-2 is currently operational and is the longest-lived lunar rover to date.
TENACIOUS
The Hakuto-R Mission 2 includes a rover called "TENACIOUS", designed and manufactured in Luxembourg which will explore the area around the landing site, after being lowered to the lunar surface from the lander.
Planned missions
Proposed missions
ATHLETE
NASA's plans for future Moon missions call for rovers that have a far longer range than the Apollo rovers. The All-Terrain Hex-Legged Extra-Terrestrial Explorer (ATHLETE) is a six-legged robotic lunar rover test-bed under development by the Jet Propulsion Laboratory (JPL). ATHLETE is a testbed for systems and is designed for use on the Moon. The system is in development along with NASA's Johnson and Ames Centers, Stanford University and Boeing. ATHLETE is designed, for maximum efficiency, to be able to both roll and walk over a wide range of terrains.
Lunar Polar Exploration Mission rover
The Lunar Polar Exploration Mission is a robotic lunar mission concept by Indian Space Research Organisation and the Japan Aerospace Exploration Agency that would send a lunar rover and lander to explore the south pole region of the Moon in 2028. The Japanese agency is likely to provide the under-development H3 launch vehicle and the rover, while the Indian agency would be responsible for the lander.
Cancelled
Lunokhod 3
Lunokhod 3 was built for a Moon landing in 1977 as Luna 25 but never flew to the Moon due to lack of launchers and funding. It remains at the NPO Lavochkin museum.
Apollo Lunar Roving Vehicle 4, 5 and 6
They would have been for Apollo 18, 19 and 20. Only the rover for Apollo 18 (LRV-4) was built. After the cancellation of that mission, it was used as spare parts for the previous rovers.
Resource Prospector
Resource Prospector is a cancelled mission concept by NASA of a rover that would have performed a survey expedition on a polar region of the Moon. The rover was to attempt to detect and map the location of volatiles such as hydrogen, oxygen and lunar water which could foster more affordable and sustainable human exploration to the Moon, Mars, and other Solar System bodies.
The mission concept was still in its pre-formulation stage when it was scrapped in April 2018. The Resource Prospector mission was proposed to be launched in 2022. Its science instruments will be flown on several commercial lander missions contracted with NASA's new Commercial Lunar Payload Services program.
VIPER
| Technology | Rovers | null |
35541901 | https://en.wikipedia.org/wiki/Thin%20disk | Thin disk | The thin disk is a structural component of spiral and S0-type galaxies, composed of stars, gas and dust. It is the main non-centre (e.g. galactic bulge) density, of such matter. That of the Milky Way is thought to have a scale height of around in the vertical axis perpendicular to the disk, and a scale length of around in the horizontal axis, in the direction of the radius. For comparison, the Sun is out from the center. The thin disk contributes about 85% of the stars in the Galactic plane and 95% of the total disk stars. It can be set apart from the thick disk of a galaxy since the latter is composed of older population stars created at an earlier stage of the galaxy formation and thus has fewer heavy elements. Stars in the thin disk, on the other hand, are created as a result of gas accretion at the later stages of a galaxy formation and are on average more metal-rich.
The thin disk contains stars with a wide range of ages and may be divided into a series of sub-populations of increasing age. Notwithstanding, it is considered to be considerably younger than the thick disk.
Based upon the emerging science of nucleocosmochronology, the Galactic thin disk of the Milky Way is estimated to have been formed 8.8 ± 1.7 billion years ago. It may have collided with a smaller satellite galaxy, causing the stars in the thin disk to be shaken up and creating the thick disk, while the gas would have settled into the galactic plane and reformed the thin disk.
| Physical sciences | Basics_2 | Astronomy |
31367277 | https://en.wikipedia.org/wiki/Grain%20size | Grain size | Grain size (or particle size) is the diameter of individual grains of sediment, or the lithified particles in clastic rocks. The term may also be applied to other granular materials. This is different from the crystallite size, which refers to the size of a single crystal inside a particle or grain. A single grain can be composed of several crystals. Granular material can range from very small colloidal particles, through clay, silt, sand, gravel, and cobbles, to boulders.
Krumbein phi scale
Size ranges define limits of classes that are given names in the Wentworth scale (or Udden–Wentworth scale named after geologists Chester K. Wentworth and Johan A. Udden) used in the United States. The Krumbein phi (φ) scale, a modification of the Wentworth scale created by W. C. Krumbein in 1934, is a logarithmic scale computed by the equation
where
is the Krumbein phi scale,
is the diameter of the particle or grain in millimeters (Krumbein and Monk's equation) and
is a reference diameter, equal to 1 mm (to make the equation dimensionally consistent).
This equation can be rearranged to find diameter using φ:
In some schemes, gravel is anything larger than sand (comprising granule, pebble, cobble, and boulder in the table above).
International scale
ISO 14688-1:2017, establishes the basic principles for identifying and classifying soils based on those material and mass characteristics most commonly used for soils for engineering purposes. ISO 14688-1 applies to natural soils in situ, similar man-made materials in situ and soils redeposited by people.
Sorting
An accumulation of sediment can also be characterized by the grain size distribution. A sediment deposit can undergo sorting when a particle size range is removed by an agency such as a river or the wind. The sorting can be quantified using the Inclusive Graphic Standard Deviation:
where
is the Inclusive Graphic Standard Deviation in phi units
is the 84th percentile of the grain size distribution in phi units, etc.
The result of this can be described using the following terms:
| Physical sciences | Sedimentology | Earth science |
20203982 | https://en.wikipedia.org/wiki/Spinneret%20%28polymers%29 | Spinneret (polymers) | A spinneret is a device used to extrude a polymer solution or polymer melt to form fibers. Streams of viscous polymer exit via the spinneret into air or liquid leading to a phase inversion which allows the polymer to solidify. The individual polymer chains tend to align in the fiber because of viscous flow. This airstream liquid-to-fiber formation process is similar to the production process for cotton candy. The fiber production process is generally referred to as "spinning". Depending on the type of spinneret used, either solid or hollow fibers can be formed.
Spinnerets are also used for electrospinning and electrospraying applications. They are sometimes called coaxial needles, or coaxial emitters.
Spinnerets are usually made of metals with melting points too low to withstand the heating processes employed in industrial metallurgy, and thus are generally not used to form metallic fibers.
| Technology | Spinning | null |
32861742 | https://en.wikipedia.org/wiki/Health%20and%20usage%20monitoring%20systems | Health and usage monitoring systems | Health and usage monitoring systems (HUMS) is a generic term given to activities that utilize data collection and analysis techniques to help ensure availability, reliability and safety of vehicles. Activities similar to, or sometimes used interchangeably with, HUMS include condition-based maintenance (CBM) and operational data recording (ODR). This term HUMS is often used in reference to airborne craft and in particular rotor-craft – the term is cited as being introduced by the offshore oil industry after a commercial Chinook crashed in the North Sea, killing all but one passenger and one crew member in 1986.
HUMS technology and regulation continues to be developed.
HUMS are now used not only for safety but for a number of other reasons including
Maintenance: reduced mission aborts, fewer instances of aircraft on ground (AOG), simplified logistics for fleet deployment
Cost: “maintain as you fly” maintenance flights are not required. Performing repairs when the damage is minor increases the aircraft mean time before failure (MTBF) and decreases the mean time to repair (MTTR)
Operational: improved flight safety, mission reliability and effectiveness
Performance: improved aircraft performance and reduced fuel consumption
Recent advances in the technology include predictive algorithms providing Remaining Useful Life estimates of components and automated wireless data transfer from the aircraft via WiFi or Cellular.
| Technology | Aircraft components | null |
2792729 | https://en.wikipedia.org/wiki/Crop%20yield | Crop yield | In agriculture, the yield is a measurement of the amount of a crop grown, or product such as wool, meat or milk produced, per unit area of land. The seed ratio is another way of calculating yields.
Innovations, such as the use of fertilizer, the creation of better farming tools, new methods of farming and improved crop varieties, have improved yields. The higher the yield and more intensive use of the farmland, the higher the productivity and profitability of a farm; this increases the well-being of farming families. Surplus crops beyond the needs of subsistence agriculture can be sold or bartered. The more grain or fodder a farmer can produce, the more draft animals such as horses and oxen could be supported and harnessed for labour and production of manure. Increased crop yields also means fewer hands are needed on farm, freeing them for industry and commerce. This, in turn, led to the formation and growth of cities, which then translated into an increased demand for foodstuffs or other agricultural products.
Measurement
The units by which the yield of a crop is usually measured today are kilograms per hectare or bushels per acre.
Long-term cereal yields in the United Kingdom were some 500 kg/ha in Medieval times, jumping to 2000 kg/ha in the Industrial Revolution, and jumping again to 8000 kg/ha in the Green Revolution. Each technological advance increasing the crop yield also reduces the society's ecological footprint.
Yields are related to agricultural productivity, but are not synonymous. Agricultural productivity is measured in money produced per unit of land, but yields are measured in the weight of the crop produced per unit of land. A farmer can invest a large amount of money to increase his yields by a few percent, for example with an extremely expensive fertilizer, but if that cost is so high that it does not produce a comparative return on investment, his profits decline, and the higher yield can mean a lower agricultural productivity in this case. A yield is a 'partial measure of productivity', because it may fail to accurately measure the actual productivity of the farming operation by not including the totality of the inputs.
Seed multiplication ratio
The seed multiplication ratio is the ratio between the investment in seed versus the yield. For example, if three grains are harvested for each grain seeded, the resulting multiplication ratio is 1:3, which is considered by some agronomists as the minimum required to sustain human life. One of the three seeds must be set aside for the next planting season, the remaining two either consumed by the grower, or for livestock feed. In parts of Europe the seed ratio during the 9th century was merely 1:2.5, in the Low Countries it improved to 1:14 with the introduction of the three-field system of crop rotation around the 14th century.
Seed multiplication ratio is variable, subject to several factors. Agricultural improvements can raise the ratio, and revisions were recommended in 2018 by the Indian Council of Agricultural Research.
Law of physiological relations
Alexander Mitscherlich studied crop yields in 1909 and articulated a "law of physiological relations". It was compared to the law of diminishing returns in 1942, when Liebig's law of the minimum and the limiting factors of Frederick Blackman were also noted:
Liebig's Law of the Minimum was the formulation of an idea that yield of a crop was determined primarily by the amounts of plant food that were present in minimum quantities. His idea was discussed later as the Limiting Factor by BLACKMAN and again by MITSCHERLICH as the Law of Physiological Relations. The latter was expressed as a logarithmic function between yield and the quantity of plant food constituents, which is virtually the Law of Diminishing Returns.
The relation was reviewed by Hans Schneeberger in 2009.
| Technology | Basics_3 | null |
2796131 | https://en.wikipedia.org/wiki/Introduction%20to%20quantum%20mechanics | Introduction to quantum mechanics | Quantum mechanics is the study of matter and its interactions with energy on the scale of atomic and subatomic particles. By contrast, classical physics explains matter and energy only on a scale familiar to human experience, including the behavior of astronomical bodies such as the moon. Classical physics is still used in much of modern science and technology. However, towards the end of the 19th century, scientists discovered phenomena in both the large (macro) and the small (micro) worlds that classical physics could not explain. The desire to resolve inconsistencies between observed phenomena and classical theory led to a revolution in physics, a shift in the original scientific paradigm: the development of quantum mechanics.
Many aspects of quantum mechanics are counterintuitive and can seem paradoxical because they describe behavior quite different from that seen at larger scales. In the words of quantum physicist Richard Feynman, quantum mechanics deals with "nature as She is—absurd". Features of quantum mechanics often defy simple explanations in everyday language. One example of this is the uncertainty principle: precise measurements of position cannot be combined with precise measurements of velocity. Another example is entanglement: a measurement made on one particle (such as an electron that is measured to have spin 'up') will correlate with a measurement on a second particle (an electron will be found to have spin 'down') if the two particles have a shared history. This will apply even if it is impossible for the result of the first measurement to have been transmitted to the second particle before the second measurement takes place.
Quantum mechanics helps us understand chemistry, because it explains how atoms interact with each other and form molecules. Many remarkable phenomena can be explained using quantum mechanics, like superfluidity. For example, if liquid helium cooled to a temperature near absolute zero is placed in a container, it spontaneously flows up and over the rim of its container; this is an effect which cannot be explained by classical physics.
History
James C. Maxwell's unification of the equations governing electricity, magnetism, and light in the late 19th century led to experiments on the interaction of light and matter. Some of these experiments had aspects which could not be explained until quantum mechanics emerged in the early part of the 20th century.
Evidence of quanta from the photoelectric effect
The seeds of the quantum revolution appear in the discovery by J.J. Thomson in 1897 that cathode rays were not continuous but "corpuscles" (electrons). Electrons had been named just six years earlier as part of the emerging theory of atoms. In 1900, Max Planck, unconvinced by the atomic theory, discovered that he needed discrete entities like atoms or electrons to explain black-body radiation.
Very hot – red hot or white hot – objects look similar when heated to the same temperature. This look results from a common curve of light intensity at different frequencies (colors), which is called black-body radiation. White hot objects have intensity across many colors in the visible range. The lowest frequencies above visible colors are infrared light, which also give off heat. Continuous wave theories of light and matter cannot explain the black-body radiation curve. Planck spread the heat energy among individual "oscillators" of an undefined character but with discrete energy capacity; this model explained black-body radiation.
At the time, electrons, atoms, and discrete oscillators were all exotic ideas to explain exotic phenomena. But in 1905 Albert Einstein proposed that light was also corpuscular, consisting of "energy quanta", in contradiction to the established science of light as a continuous wave, stretching back a hundred years to Thomas Young's work on diffraction.
Einstein's revolutionary proposal started by reanalyzing Planck's black-body theory, arriving at the same conclusions by using the new "energy quanta". Einstein then showed how energy quanta connected to Thomson's electron. In 1902, Philipp Lenard directed light from an arc lamp onto freshly cleaned metal plates housed in an evacuated glass tube. He measured the electric current coming off the metal plate, at higher and lower intensities of light and for different metals. Lenard showed that amount of current – the number of electrons – depended on the intensity of the light, but that the velocity of these electrons did not depend on intensity. This is the photoelectric effect. The continuous wave theories of the time predicted that more light intensity would accelerate the same amount of current to higher velocity, contrary to this experiment. Einstein's energy quanta explained the volume increase: one electron is ejected for each quantum: more quanta mean more electrons.
Einstein then predicted that the electron velocity would increase in direct proportion to the light frequency above a fixed value that depended upon the metal. Here the idea is that energy in energy-quanta depends upon the light frequency; the energy transferred to the electron comes in proportion to the light frequency. The type of metal gives a barrier, the fixed value, that the electrons must climb over to exit their atoms, to be emitted from the metal surface and be measured.
Ten years elapsed before Millikan's definitive experiment verified Einstein's prediction. During that time many scientists rejected the revolutionary idea of quanta. But Planck's and Einstein's concept was in the air and soon began to affect other physics and quantum theories.
Quantization of bound electrons in atoms
Experiments with light and matter in the late 1800s uncovered a reproducible but puzzling regularity. When light was shown through purified gases, certain frequencies (colors) did not pass. These dark absorption 'lines' followed a distinctive pattern: the gaps between the lines decreased steadily. By 1889, the Rydberg formula predicted the lines for hydrogen gas using only a constant number and the integers to index the lines. The origin of this regularity was unknown. Solving this mystery would eventually become the first major step toward quantum mechanics.
Throughout the 19th century evidence grew for the atomic nature of matter. With Thomson's discovery of the electron in 1897, scientist began the search for a model of the interior of the atom. Thomson proposed negative electrons swimming in a pool of positive charge. Between 1908 and 1911, Rutherford showed that the positive part was only 1/3000th of the diameter of the atom.
Models of "planetary" electrons orbiting a nuclear "Sun" were proposed, but cannot explain why the electron does not simply fall into the positive charge. In 1913 Niels Bohr and Ernest Rutherford connected the new atom models to the mystery of the Rydberg formula: the orbital radius of the electrons were constrained and the resulting energy differences matched the energy differences in the absorption lines. This meant that absorption and emission of light from atoms was energy quantized: only specific energies that matched the difference in orbital energy would be emitted or absorbed.
Trading one mystery – the regular pattern of the Rydberg formula – for another mystery – constraints on electron orbits – might not seem like a big advance, but the new atom model summarized many other experimental findings. The quantization of the photoelectric effect and now the quantization of the electron orbits set the stage for the final revolution.
Throughout the first and the modern era of quantum mechanics the concept that classical mechanics must be valid macroscopically constrained possible quantum models. This concept was formalized by Bohr in 1923 as the correspondence principle. It requires quantum theory to converge to classical limits.
A related concept is Ehrenfest's theorem, which shows that the average values obtained from quantum mechanics (e.g. position and momentum) obey classical laws.
Quantization of spin
In 1922 Otto Stern and Walther Gerlach demonstrated that the magnetic properties of silver atoms defy classical explanation, the work contributing to Stern’s 1943 Nobel Prize in Physics. They fired a beam of silver atoms through a magnetic field. According to classical physics, the atoms should have emerged in a spray, with a continuous range of directions. Instead, the beam separated into two, and only two, diverging streams of atoms. Unlike the other quantum effects known at the time, this striking result involves the state of a single atom. In 1927, T.E. Phipps and J.B. Taylor obtained a similar, but less pronounced effect using hydrogen atoms in their ground state, thereby eliminating any doubts that may have been caused by the use of silver atoms.
In 1924, Wolfgang Pauli called it "two-valuedness not describable classically" and associated it with electrons in the outermost shell. The experiments lead to formulation of its theory described to arise from spin of the electron in 1925, by Samuel Goudsmit and George Uhlenbeck, under the advice of Paul Ehrenfest.
Quantization of matter
In 1924 Louis de Broglie proposed that electrons in an atom are constrained not in "orbits" but as standing waves. In detail his solution did not work, but his hypothesis – that the electron "corpuscle" moves in the atom as a wave – spurred Erwin Schrödinger to develop a wave equation for electrons; when applied to hydrogen the Rydberg formula was accurately reproduced.
Max Born's 1924 paper "Zur Quantenmechanik" was the first use of the words "quantum mechanics" in print. His later work included developing quantum collision models; in a footnote to a 1926 paper he proposed the Born rule connecting theoretical models to experiment.
In 1927 at Bell Labs, Clinton Davisson and Lester Germer fired slow-moving electrons at a crystalline nickel target which showed a diffraction pattern indicating wave nature of electron whose theory was fully explained by Hans Bethe. A similar experiment by George Paget Thomson and Alexander Reid, firing electrons at thin celluloid foils and later metal films, observing rings, independently discovered matter wave nature of electrons.
Further developments
In 1928 Paul Dirac published his relativistic wave equation simultaneously incorporating relativity, predicting anti-matter, and providing a complete theory for the Stern–Gerlach result. These successes launched a new fundamental understanding of our world at small scale: quantum mechanics.
Planck and Einstein started the revolution with quanta that broke down the continuous models of matter and light. Twenty years later "corpuscles" like electrons came to be modeled as continuous waves. This result came to be called wave-particle duality, one iconic idea along with the uncertainty principle that sets quantum mechanics apart from older models of physics.
Quantum radiation, quantum fields
In 1923 Compton demonstrated that the Planck-Einstein energy quanta from light also had momentum; three years later the "energy quanta" got a new name "photon" Despite its role in almost all stages of the quantum revolution, no explicit model for light quanta existed until 1927 when Paul Dirac began work on a quantum theory of radiation that became quantum electrodynamics. Over the following decades this work evolved into quantum field theory, the basis for modern quantum optics and particle physics.
Wave–particle duality
The concept of wave–particle duality says that neither the classical concept of "particle" nor of "wave" can fully describe the behavior of quantum-scale objects, either photons or matter. Wave–particle duality is an example of the principle of complementarity in quantum physics. An elegant example of wave-particle duality is the double-slit experiment.
In the double-slit experiment, as originally performed by Thomas Young in 1803, and then Augustin Fresnel a decade later, a beam of light is directed through two narrow, closely spaced slits, producing an interference pattern of light and dark bands on a screen. The same behavior can be demonstrated in water waves: the double-slit experiment was seen as a demonstration of the wave nature of light.
Variations of the double-slit experiment have been performed using electrons, atoms, and even large molecules, and the same type of interference pattern is seen. Thus it has been demonstrated that all matter possesses wave characteristics.
If the source intensity is turned down, the same interference pattern will slowly build up, one "count" or particle (e.g. photon or electron) at a time. The quantum system acts as a wave when passing through the double slits, but as a particle when it is detected. This is a typical feature of quantum complementarity: a quantum system acts as a wave in an experiment to measure its wave-like properties, and like a particle in an experiment to measure its particle-like properties. The point on the detector screen where any individual particle shows up is the result of a random process. However, the distribution pattern of many individual particles mimics the diffraction pattern produced by waves.
Uncertainty principle
Suppose it is desired to measure the position and speed of an object—for example, a car going through a radar speed trap. It can be assumed that the car has a definite position and speed at a particular moment in time. How accurately these values can be measured depends on the quality of the measuring equipment. If the precision of the measuring equipment is improved, it provides a result closer to the true value. It might be assumed that the speed of the car and its position could be operationally defined and measured simultaneously, as precisely as might be desired.
In 1927, Heisenberg proved that this last assumption is not correct. Quantum mechanics shows that certain pairs of physical properties, for example, position and speed, cannot be simultaneously measured, nor defined in operational terms, to arbitrary precision: the more precisely one property is measured, or defined in operational terms, the less precisely can the other be thus treated. This statement is known as the uncertainty principle. The uncertainty principle is not only a statement about the accuracy of our measuring equipment but, more deeply, is about the conceptual nature of the measured quantities—the assumption that the car had simultaneously defined position and speed does not work in quantum mechanics. On a scale of cars and people, these uncertainties are negligible, but when dealing with atoms and electrons they become critical.
Heisenberg gave, as an illustration, the measurement of the position and momentum of an electron using a photon of light. In measuring the electron's position, the higher the frequency of the photon, the more accurate is the measurement of the position of the impact of the photon with the electron, but the greater is the disturbance of the electron. This is because from the impact with the photon, the electron absorbs a random amount of energy, rendering the measurement obtained of its momentum increasingly uncertain, for one is necessarily measuring its post-impact disturbed momentum from the collision products and not its original momentum (momentum which should be simultaneously measured with position). With a photon of lower frequency, the disturbance (and hence uncertainty) in the momentum is less, but so is the accuracy of the measurement of the position of the impact.
At the heart of the uncertainty principle is a fact that for any mathematical analysis in the position and velocity domains, achieving a sharper (more precise) curve in the position domain can only be done at the expense of a more gradual (less precise) curve in the speed domain, and vice versa. More sharpness in the position domain requires contributions from more frequencies in the speed domain to create the narrower curve, and vice versa. It is a fundamental tradeoff inherent in any such related or complementary measurements, but is only really noticeable at the smallest (Planck) scale, near the size of elementary particles.
The uncertainty principle shows mathematically that the product of the uncertainty in the position and momentum of a particle (momentum is velocity multiplied by mass) could never be less than a certain value, and that this value is related to the Planck constant.
Wave function collapse
Wave function collapse means that a measurement has forced or converted a quantum (probabilistic or potential) state into a definite measured value. This phenomenon is only seen in quantum mechanics rather than classical mechanics.
For example, before a photon actually "shows up" on a detection screen it can be described only with a set of probabilities for where it might show up. When it does appear, for instance in the CCD of an electronic camera, the time and space where it interacted with the device are known within very tight limits. However, the photon has disappeared in the process of being captured (measured), and its quantum wave function has disappeared with it. In its place, some macroscopic physical change in the detection screen has appeared, e.g., an exposed spot in a sheet of photographic film, or a change in electric potential in some cell of a CCD.
Eigenstates and eigenvalues
Because of the uncertainty principle, statements about both the position and momentum of particles can assign only a probability that the position or momentum has some numerical value. Therefore, it is necessary to formulate clearly the difference between the state of something indeterminate, such as an electron in a probability cloud, and the state of something having a definite value. When an object can definitely be "pinned-down" in some respect, it is said to possess an eigenstate.
In the Stern–Gerlach experiment discussed above, the quantum model predicts two possible values of spin for the atom compared to the magnetic axis. These two eigenstates are named arbitrarily 'up' and 'down'. The quantum model predicts these states will be measured with equal probability, but no intermediate values will be seen. This is what the Stern–Gerlach experiment shows.
The eigenstates of spin about the vertical axis are not simultaneously eigenstates of spin about the horizontal axis, so this atom has an equal probability of being found to have either value of spin about the horizontal axis. As described in the section above, measuring the spin about the horizontal axis can allow an atom that was spun up to spin down: measuring its spin about the horizontal axis collapses its wave function into one of the eigenstates of this measurement, which means it is no longer in an eigenstate of spin about the vertical axis, so can take either value.
The Pauli exclusion principle
In 1924, Wolfgang Pauli proposed a new quantum degree of freedom (or quantum number), with two possible values, to resolve inconsistencies between observed molecular spectra and the predictions of quantum mechanics. In particular, the spectrum of atomic hydrogen had a doublet, or pair of lines differing by a small amount, where only one line was expected. Pauli formulated his exclusion principle, stating, "There cannot exist an atom in such a quantum state that two electrons within [it] have the same set of quantum numbers."
A year later, Uhlenbeck and Goudsmit identified Pauli's new degree of freedom with the property called spin whose effects were observed in the Stern–Gerlach experiment.
Dirac wave equation
In 1928, Paul Dirac extended the Pauli equation, which described spinning electrons, to account for special relativity. The result was a theory that dealt properly with events, such as the speed at which an electron orbits the nucleus, occurring at a substantial fraction of the speed of light. By using the simplest electromagnetic interaction, Dirac was able to predict the value of the magnetic moment associated with the electron's spin and found the experimentally observed value, which was too large to be that of a spinning charged sphere governed by classical physics. He was able to solve for the spectral lines of the hydrogen atom and to reproduce from physical first principles Sommerfeld's successful formula for the fine structure of the hydrogen spectrum.
Dirac's equations sometimes yielded a negative value for energy, for which he proposed a novel solution: he posited the existence of an antielectron and a dynamical vacuum. This led to the many-particle quantum field theory.
Quantum entanglement
In quantum physics, a group of particles can interact or be created together in such a way that the quantum state of each particle of the group cannot be described independently of the state of the others, including when the particles are separated by a large distance. This is known as quantum entanglement.
An early landmark in the study of entanglement was the Einstein–Podolsky–Rosen (EPR) paradox, a thought experiment proposed by Albert Einstein, Boris Podolsky and Nathan Rosen which argues that the description of physical reality provided by quantum mechanics is incomplete. In a 1935 paper titled "Can Quantum-Mechanical Description of Physical Reality be Considered Complete?", they argued for the existence of "elements of reality" that were not part of quantum theory, and speculated that it should be possible to construct a theory containing these hidden variables.
The thought experiment involves a pair of particles prepared in what would later become known as an entangled state. Einstein, Podolsky, and Rosen pointed out that, in this state, if the position of the first particle were measured, the result of measuring the position of the second particle could be predicted. If instead the momentum of the first particle were measured, then the result of measuring the momentum of the second particle could be predicted. They argued that no action taken on the first particle could instantaneously affect the other, since this would involve information being transmitted faster than light, which is forbidden by the theory of relativity. They invoked a principle, later known as the "EPR criterion of reality", positing that: "If, without in any way disturbing a system, we can predict with certainty (i.e., with probability equal to unity) the value of a physical quantity, then there exists an element of reality corresponding to that quantity." From this, they inferred that the second particle must have a definite value of both position and of momentum prior to either quantity being measured. But quantum mechanics considers these two observables incompatible and thus does not associate simultaneous values for both to any system. Einstein, Podolsky, and Rosen therefore concluded that quantum theory does not provide a complete description of reality. In the same year, Erwin Schrödinger used the word "entanglement" and declared: "I would not call that one but rather the characteristic trait of quantum mechanics."
The Irish physicist John Stewart Bell carried the analysis of quantum entanglement much further. He deduced that if measurements are performed independently on the two separated particles of an entangled pair, then the assumption that the outcomes depend upon hidden variables within each half implies a mathematical constraint on how the outcomes on the two measurements are correlated. This constraint would later be named the Bell inequality. Bell then showed that quantum physics predicts correlations that violate this inequality. Consequently, the only way that hidden variables could explain the predictions of quantum physics is if they are "nonlocal", which is to say that somehow the two particles are able to interact instantaneously no matter how widely they ever become separated. Performing experiments like those that Bell suggested, physicists have found that nature obeys quantum mechanics and violates Bell inequalities. In other words, the results of these experiments are incompatible with any local hidden variable theory.
Quantum field theory
The idea of quantum field theory began in the late 1920s with British physicist Paul Dirac, when he attempted to quantize the energy of the electromagnetic field; just as in quantum mechanics the energy of an electron in the hydrogen atom was quantized. Quantization is a procedure for constructing a quantum theory starting from a classical theory.
Merriam-Webster defines a field in physics as "a region or space in which a given effect (such as magnetism) exists". Other effects that manifest themselves as fields are gravitation and static electricity. In 2008, physicist Richard Hammond wrote:
Sometimes we distinguish between quantum mechanics (QM) and quantum field theory (QFT). QM refers to a system in which the number of particles is fixed, and the fields (such as the electromechanical field) are continuous classical entities. QFT ... goes a step further and allows for the creation and annihilation of particles ...
He added, however, that quantum mechanics is often used to refer to "the entire notion of quantum view".
In 1931, Dirac proposed the existence of particles that later became known as antimatter. Dirac shared the Nobel Prize in Physics for 1933 with Schrödinger "for the discovery of new productive forms of atomic theory".
Quantum electrodynamics
Quantum electrodynamics (QED) is the name of the quantum theory of the electromagnetic force. Understanding QED begins with understanding electromagnetism. Electromagnetism can be called "electrodynamics" because it is a dynamic interaction between electrical and magnetic forces. Electromagnetism begins with the electric charge.
Electric charges are the sources of and create, electric fields. An electric field is a field that exerts a force on any particles that carry electric charges, at any point in space. This includes the electron, proton, and even quarks, among others. As a force is exerted, electric charges move, a current flows, and a magnetic field is produced. The changing magnetic field, in turn, causes electric current (often moving electrons). The physical description of interacting charged particles, electrical currents, electrical fields, and magnetic fields is called electromagnetism.
In 1928 Paul Dirac produced a relativistic quantum theory of electromagnetism. This was the progenitor to modern quantum electrodynamics, in that it had essential ingredients of the modern theory. However, the problem of unsolvable infinities developed in this relativistic quantum theory. Years later, renormalization largely solved this problem. Initially viewed as a provisional, suspect procedure by some of its originators, renormalization eventually was embraced as an important and self-consistent tool in QED and other fields of physics. Also, in the late 1940s Feynman diagrams provided a way to make predictions with QED by finding a probability amplitude for each possible way that an interaction could occur. The diagrams showed in particular that the electromagnetic force is the exchange of photons between interacting particles.
The Lamb shift is an example of a quantum electrodynamics prediction that has been experimentally verified. It is an effect whereby the quantum nature of the electromagnetic field makes the energy levels in an atom or ion deviate slightly from what they would otherwise be. As a result, spectral lines may shift or split.
Similarly, within a freely propagating electromagnetic wave, the current can also be just an abstract displacement current, instead of involving charge carriers. In QED, its full description makes essential use of short-lived virtual particles. There, QED again validates an earlier, rather mysterious concept.
Standard Model
The Standard Model of particle physics is the quantum field theory that describes three of the four known fundamental forces (electromagnetic, weak and strong interactions – excluding gravity) in the universe and classifies all known elementary particles. It was developed in stages throughout the latter half of the 20th century, through the work of many scientists worldwide, with the current formulation being finalized in the mid-1970s upon experimental confirmation of the existence of quarks. Since then, proof of the top quark (1995), the tau neutrino (2000), and the Higgs boson (2012) have added further credence to the Standard Model. In addition, the Standard Model has predicted various properties of weak neutral currents and the W and Z bosons with great accuracy.
Although the Standard Model is believed to be theoretically self-consistent and has demonstrated success in providing experimental predictions, it leaves some physical phenomena unexplained and so falls short of being a complete theory of fundamental interactions. For example, it does not fully explain baryon asymmetry, incorporate the full theory of gravitation as described by general relativity, or account for the universe's accelerating expansion as possibly described by dark energy. The model does not contain any viable dark matter particle that possesses all of the required properties deduced from observational cosmology. It also does not incorporate neutrino oscillations and their non-zero masses. Accordingly, it is used as a basis for building more exotic models that incorporate hypothetical particles, extra dimensions, and elaborate symmetries (such as supersymmetry) to explain experimental results at variance with the Standard Model, such as the existence of dark matter and neutrino oscillations.
Interpretations
The physical measurements, equations, and predictions pertinent to quantum mechanics are all consistent and hold a very high level of confirmation. However, the question of what these abstract models say about the underlying nature of the real world has received competing answers. These interpretations are widely varying and sometimes somewhat abstract. For instance, the Copenhagen interpretation states that before a measurement, statements about a particle's properties are completely meaningless, while the many-worlds interpretation describes the existence of a multiverse made up of every possible universe.
Light behaves in some aspects like particles and in other aspects like waves. Matter—the "stuff" of the universe consisting of particles such as electrons and atoms—exhibits wavelike behavior too. Some light sources, such as neon lights, give off only certain specific frequencies of light, a small set of distinct pure colors determined by neon's atomic structure. Quantum mechanics shows that light, along with all other forms of electromagnetic radiation, comes in discrete units, called photons, and predicts its spectral energies (corresponding to pure colors), and the intensities of its light beams. A single photon is a quantum, or smallest observable particle, of the electromagnetic field. A partial photon is never experimentally observed. More broadly, quantum mechanics shows that many properties of objects, such as position, speed, and angular momentum, that appeared continuous in the zoomed-out view of classical mechanics, turn out to be (in the very tiny, zoomed-in scale of quantum mechanics) quantized. Such properties of elementary particles are required to take on one of a set of small, discrete allowable values, and since the gap between these values is also small, the discontinuities are only apparent at very tiny (atomic) scales.
Applications
Everyday applications
The relationship between the frequency of electromagnetic radiation and the energy of each photon is why ultraviolet light can cause sunburn, but visible or infrared light cannot. A photon of ultraviolet light delivers a high amount of energy—enough to contribute to cellular damage such as occurs in a sunburn. A photon of infrared light delivers less energy—only enough to warm one's skin. So, an infrared lamp can warm a large surface, perhaps large enough to keep people comfortable in a cold room, but it cannot give anyone a sunburn.
Technological applications
Applications of quantum mechanics include the laser, the transistor, the electron microscope, and magnetic resonance imaging. A special class of quantum mechanical applications is related to macroscopic quantum phenomena such as superfluid helium and superconductors. The study of semiconductors led to the invention of the diode and the transistor, which are indispensable for modern electronics.
In even a simple light switch, quantum tunneling is absolutely vital, as otherwise the electrons in the electric current could not penetrate the potential barrier made up of a layer of oxide. Flash memory chips found in USB drives also use quantum tunneling, to erase their memory cells.
| Physical sciences | Quantum mechanics | Physics |
2796768 | https://en.wikipedia.org/wiki/Rose%20%28color%29 | Rose (color) | Rose is the color halfway between red and magenta on the HSV color wheel, also known as the RGB color wheel, on which it is at hue angle of 330 degrees.
Rose is one of the tertiary colors on the HSV (RGB) color wheel. The complementary color of rose is spring green. Sometimes rose is quoted instead as the web-safe color FF00CC, which is closer to magenta than to red, corresponding to a hue angle near 320 degrees, or the web-safe color FF0077, which is closer to red than magenta, corresponding to a hue angle of about 340 degrees.
Shades of rose
Etymology of rose
The first recorded use of rose as a color name in English was in 1382.
The etymology of the color name rose is the same as that of the name of the rose flower. The name originates from Latin rosa, borrowed through Oscan from colonial Greek in southern Italy: rhodon (Aeolic form: wrodon), from Aramaic wurrdā, from Assyrian wurtinnu, from Old Iranian *warda (cf. Avestan warda, Sogdian ward, Parthian wâr).
In culture
Geography
British historian John William Burgon famously described the Jordanian city of Petra as being colored rose, writing:
Match me such marvel save in Eastern clime,
A rose-red city – half as old as time!
Marrakech, Morocco is called the Rose City because many of its buildings are colored various tones of rose.
Portland, Oregon is nicknamed "The Rose City" for the number of roses and rose gardens that thrive there.
Music
"La Vie en rose" (French for "Life through rose-coloured glasses", literally "Life in pink") was the signature song of French singer Édith Piaf. Piaf first popularized the song in 1946. It has been covered by many artists since.
Rose Colored Glasses is the 1978 debut album by country singer-songwriter John Conlee.
Occult
According to New Age author C. W. Leadbeater, who claimed to be clairvoyant, of the seven types of etheric atoms that he claimed to be able to observe with his third eye circulating through the human etheric body (colored violet, blue, green, yellow, orange, dark red, and rose), the flow of the rose colored etheric atoms (also called by Leadbeater the rose vitality globule) from the sun into the rainbow colored spleen chakra is the most important since all the other etheric atoms are derived from it and the rose colored atom vivifies the nervous system. Leadbeater also asserted that humans feel good around pine trees because they radiate more rose colored etheric atoms than any other plant.
Politics
The revolution in which previous Georgian president Mikhail Saakashvili came to power in 2003 was called the "Rose Revolution"
The color is used by the French Socialist Party
Religion
In the Latin liturgical rites of the Catholic Church, priests may wear rose colored vestments on Gaudete Sunday (the third Sunday of Advent) and Laetare Sunday (the fourth Sunday of Lent).
| Physical sciences | Colors | Physics |
3755359 | https://en.wikipedia.org/wiki/Confluence | Confluence | In geography, a confluence (also: conflux) occurs where two or more watercourses join to form a single channel. A confluence can occur in several configurations: at the point where a tributary joins a larger river (main stem); or where two streams meet to become the source of a river of a new name (such as the confluence of the Monongahela and Allegheny rivers, forming the Ohio River); or where two separated channels of a river (forming a river island) rejoin at the downstream end. The point of confluence where the channel flows into a larger body of water may be called the river mouth.
Scientific study of confluences
Confluences are studied in a variety of sciences. Hydrology studies the characteristic flow patterns of confluences and how they give rise to patterns of erosion, bars, and scour pools. The water flows and their consequences are often studied with mathematical models. Confluences are relevant to the distribution of living organisms (i.e., ecology) as well; "the general pattern [downstream of confluences] of increasing stream flow and decreasing slopes drives a corresponding shift in habitat characteristics."
Another science relevant to the study of confluences is chemistry, because sometimes the mixing of the waters of two streams triggers a chemical reaction, particularly in a polluted stream. The United States Geological Survey gives an example: "chemical changes occur when a stream contaminated with acid mine drainage combines with a stream with near-neutral pH water; these reactions happen very rapidly and influence the subsequent transport of metals downstream of the mixing zone."
A natural phenomenon at confluences that is obvious even to casual observers is a difference in color between the two streams; see images in this article for several examples. According to Lynch, "the color of each river is determined by many things: type and amount of vegetation in the watershed, geological properties, dissolved chemicals, sediments and biologic content – usually algae." Lynch also notes that color differences can persist for miles downstream before they finally blend completely.
River confluence flow zones
Hydrodynamic behaviour of flow in a confluence can be divided into six distinct features which are commonly called confluence flow zones (CFZ). These include
Stagnation zone
Flow deflection zone
Flow separation zone / recirculation zone
Maximum velocity zone
Flow recovery zone
Shear layers
Confluences in engineering
The broader field of engineering encompasses a vast assortment of subjects which concern confluences.
In hydraulic civil engineering, where two or more underground culverted / artificially buried watercourses intersect, great attention should be paid to the hydrodynamic aspects of the system to ensure the longevity and efficiency of the structure.
Engineers have to design these systems whilst considering a list of factors that ensure the discharge point is structurally stable as the entrance of the lateral culvert into the main structure may compromise the stability of the structure due to the lack of support at the discharge, this often constitutes additional supports in the form of structural bracing. The velocities and hydraulic efficiencies should be meticulously calculated and can be altered by integrating different combinations of geometries, components such a gradients, cascades and an adequate junction angle which is sympathetic to the direction of the watercourse’s flow to minimise turbulent flow, maximise evacuation velocity and to ultimately maximise hydraulic efficiency.
Cultural and societal significance
Since rivers often serve as political boundaries, confluences sometimes demarcate three abutting political entities, such as nations, states, or provinces, forming a tripoint. Various examples are found in the list below.
A number of major cities, such as Chongqing, St. Louis, and Khartoum, arose at confluences; further examples appear in the list. Within a city, a confluence often forms a visually prominent point, so that confluences are sometimes chosen as the site of prominent public buildings or monuments, as in Koblenz, Lyon, and Winnipeg. Cities also often build parks at confluences, sometimes as projects of municipal improvement, as at Portland and Pittsburgh. In other cases, a confluence is an industrial site, as in Philadelphia or Mannheim. Often a confluence lies in the shared floodplain of the two rivers and nothing is built on it, for example at Manaus, described below.
One other way that confluences may be exploited by humans is as sacred places in religions. Rogers suggests that for the ancient peoples of the Iron Age in northwest Europe, watery locations were often sacred, especially sources and confluences. Pre-Christian Slavic peoples chose confluences as the sites for fortified triangular temples, where they practiced human sacrifice and other sacred rites. In Hinduism, the confluence of two sacred rivers often is a pilgrimage site for ritual bathing. In Pittsburgh, a number of adherents to Mayanism consider their city's confluence to be sacred.
Notable confluences
Africa
At Lokoja, Nigeria, the Benue River flows into the Niger.
At Kazungula in Zambia, the Chobe River flows into the Zambezi. The confluence defines the tripoint of Zambia (north of the rivers), Botswana (south of the rivers) and Namibia (west of the rivers). The land border between Botswana and Zimbabwe to the east also reaches the Zambezi at this confluence, so there is a second tripoint (Zambia-Botswana-Zimbabwe) only 150 meters downstream from the first. See Kazungula and Quadripoint, and Gallery below for image.
The Sudanese capital of Khartoum is located at the confluence of the White Nile and the Blue Nile, the beginning of the Nile.
Asia
82 km north of Basra in Iraq at the town of Al-Qurnah is the confluence of the rivers Tigris and Euphrates, forming the Shatt al-Arab.
At Devprayag in India, the Ganges River originates at the confluence of the Bhagirathi and the Alaknanda; see images above.
Near Allahabad, India, the Yamuna flows into the Ganges. In Hinduism, this is a pilgrimage site for ritual bathing; during a Kumbh Mela event tens of millions of people visit the site. In Hindu belief the site is held to be a triple confluence (Triveni Sangam), the third river being the metaphysical (not physically present) Sarasvati.
Karad, in Maharashtra, India, is the site of the Pritisangam (meaning: Lovely Confluence), a T-shaped confluence of Krishna River and Koyna River, where Koyna River mergers into Krishna River forming a T-shape and then the merged rivers flow to the east as Krishna River.
Kuala Lumpur, the capital of Malaysia, is where the Gombak River (previously known as Sungai Lumpur, which means "muddy river") flows into the Klang River at the site of the Jamek Mosque. Recently, the Kolam Biru (Blue Pool), a pool with elaborate fountains, has been installed at the apex of the confluence.
Both Taipei and New Taipei are where the Dahan and Xindian meet and flow into the Tamsui River.
The Nam Khan River flows into the Mekong at Luang Prabang in Laos.
Pak Nam Pho, the downtown of Nakhon Sawan in Thailand, is the confluence of the rivers Ping and Nan, forming the Chao Phraya the main artery of central Thailand.
Pak Phraek, the old town zone of Kanchanaburi in Thailand, is the confluence of the rivers Khwae Yai and Khwae Noi, forming the Mae Klong the main artery of western Thailand.
The Pa Sak combines the Chao Phraya at Ayutthaya in Thailand. The confluence is a location of the historic monastery Wat Phanan Choeng, built around 26 years before the founding of the Ayutthaya Kingdom.
The Jialing flows into the Yangtze at Chongqing in China. The confluence forms a focal point in the city, marked by Chaotianmen Square, built in 1998.
In the Far East, the Amur forms the international boundary between China and Russia. The Ussuri, which also demarcates the border, flows into the Amur at a point midway between Fuyuan in China and Khabarovsk in Russia. The apex of the confluence is located in a rural area, part of China, where a commemorative park, Dongji Square, has been built; it features an enormous sculpture representing the Chinese character for "East". The Amur-Ussuri border region was the location of the Sino-Soviet border conflict of 1969; the borderline near the confluence was settled peacefully by treaty in 2008.
In Georgia, in the town of Pasanauri on the southern slopes of the Caucasus Mountains, the Tetri Aragvi ("White Aragvi") is joined by the Shavi Aragvi ("Black Aragvi"). Together, these two rivers continue as the Aragvi River. The conflux is known for its dramatic visual contrast of the two rivers.
Australia
The two largest rivers in Australia, the Murray and its tributary the Darling, converge at Wentworth, New South Wales.
Europe
Seine
The Seine divides in the historical center of Paris, flowing around two river islands, the Île Saint-Louis and the Île de la Cité. At the downstream confluence, where the river becomes a single channel again, the Île de la Cité is crossed by the famous Pont Neuf, adjacent to an equestrian statue of King Henri IV and the historically more recent Vert Galant park. The site has repeatedly been portrayed by artists including Monet, Renoir, and Pissarro.
Further upstream, the Marne empties into the Seine at Charenton-le-Pont and Alfortville, just southeast of the Paris city limits. The site is dominated by the Huatian Chinagora, a four-star hotel under Chinese management.
Rhine
The Rhine carries much river traffic, and major inland ports are found at its confluence with the Ruhr at Duisburg, and with the Neckar at Mannheim; see Mannheim Harbour.
The Main flows into the Rhine just south of Mainz.
The Mosel flows into the Rhine further north at Koblenz. The name "Koblenz" itself has its origin in the Latin name "Confluentes". In German, this confluence is known as the "Deutsches Eck" ("German corner") and is the site of an imposing monument to German unification featuring an equestrian statue of Kaiser Wilhelm I.
Upstream in Switzerland, a small town also named Koblenz (for the same reason) is where the Aare joins the Rhine.
Danube basin
Passau, Germany, sometimes called the (City of Three Rivers), is the site of a triple confluence, described thus in a guidebook: "from the north the little Ilz sluices brackish water down from the peat-rich Bavarian Forest, meeting the cloudy brown of the Danube as it flows from the west and the pale snow-melt jade of the Inn from the south [i.e., the Alps] to create a murky tricolour."
The Thaya flows into the Morava in a rural location near Hohenau an der March in Austria, forming the tripoint of Austria, Czechia, and Slovakia.
The Morava flows into the Danube at Devín, on the border between Slovakia and Austria.
The Sava flows into the Danube at Belgrade, the capital of Serbia.
In karst topography, which arises in soluble rock, rivers sometimes flow underground and form subterranean confluences, as at Planina Cave in Slovenia, where the Pivka and Rak merge to form the Unica.
Other
Lyon, France lies where the Saône flows into the Rhone. A major new museum of science and anthropology, the Musée des Confluences, opened on the site in 2014.
Near Toulouse, France lies where the Ariège (river) flows into the Garonne. Both take their source in the Pyrenees.
The Lusatian Neisse flows into the Oder at a rural location in Poland opposite the German village of Ratzdorf. The two rivers form the Oder-Neisse line, the postwar boundary of Germany and Poland.
The Triangle of Three Emperors, a former political tripoint, lies in present-day Poland. The empires that abutted (in the decades before World War I) were the Austrian, German, and Russian.
Rovaniemi, the capital of Finnish Lapland and one of the largest towns above the Arctic Circle, is at the confluence of rivers Ounasjoki and Kemijoki.
Kryvyi Rih, Ukraine is located (and named after) on the confluence of the Saksahan and Inhulets River.
The Oka flows into the Volga at Nizhny Novgorod in Russia. The Alexander Nevsky Cathedral overlooks the site.
The English city of Southampton is built at the confluence of the tidal estuaries of the River Test and River Itchen which combine to form Southampton Water estuary.
North America
Mississippi basin
The Greater Twin Cities area of Minneapolis and St. Paul, Minnesota features two important Mississippi confluences. Near historical Fort Snelling and the town of Mendota—about 9 miles downstream on the Mississippi from Minneapolis—the Minnesota River flows into the Mississippi at Pike Island. The area around this confluence is a location of spiritual, cultural, and historical significance to the Dakota people and is also the site of the earliest European settlements in the Twin Cities area. About 30 miles further downstream from the Minnesota-Mississippi confluence—and 25 miles downstream from St. Paul—the Mississippi joins with the St. Croix River near Hastings, Minnesota, and Prescott, Wisconsin.
Vicksburg, Mississippi lies atop bluffs overlooking the confluence of the Mississippi River with its tributary the Yazoo. Both rivers, as well as the bluffs, played an important role in the Vicksburg Campaign, a pivotal event of the American Civil War.
The Missouri River flows into the Mississippi River at Jones-Confluence Point State Park, just north of St. Louis, Missouri. Slightly further upstream, the Illinois River flows into the Mississippi.
The Madison, Jefferson and Gallatin Rivers in Three Forks, Montana form the confluence of the Missouri River.
At Keokuk, Iowa, the Des Moines River flows into the Mississippi. This forms the political tripoint between the U.S. states of Iowa, Missouri, and Illinois.
Just south of Cairo, Illinois, the Ohio River flows into the Mississippi, forming the tripoint between the states of Illinois, Missouri, and Kentucky.
The Ohio River is formed by the confluence of the Monongahela and Allegheny rivers, located in Pittsburgh, Pennsylvania. The site is of great historical significance; in the 1970s it was upgraded by the creation of Point State Park, highlighted by a large fountain.
Atlantic watersheds
At Harpers Ferry, West Virginia, the Shenandoah River flows into the Potomac River, at the tripoint of the U.S. states of Virginia, West Virginia, and Maryland.
At Philadelphia, Pennsylvania, the Schuylkill River flows into the Delaware River, next to the former Philadelphia Naval Shipyard; the site remains industrial.
At Cohoes, New York, a few miles north of Albany, the Mohawk River flows into the Hudson in three channels separated by islands. The confluence is historically important: upstream traffic on or along the Hudson often took a left turn at the Mohawk, which offers a uniquely level passageway through the Appalachian Mountains that assisted commerce and the settlement of the West.
At Ottawa, the capital of Canada, the Rideau River flows—unusually, as a waterfall—into the Ottawa River; see Rideau Falls. On the island separating the two portions of the falls is a park with military monuments, among them the Ottawa Memorial.
The Hochelaga Archipelago, including the island and city of Montreal, is located where the Ottawa River flows into the St. Lawrence River in Quebec, Canada.
Winnipeg, Canada, is at the confluence of the Red River, and the Assiniboine River. The area is referred to as The Forks by locals, and has been an important trade location for over 6000 years.
Pacific watersheds
The Green River flows into the Colorado River at the heart of Canyonlands National Park in Utah's Canyon Country.
The Snake River flows into the Columbia River at Sacagawea State park near the Tri-Cities of Washington. It should also be noted that the significant Yakima river also flows into the Columbia just a few miles upstream, thus giving the region the unofficial preposition of Three Rivers
In Portland, Oregon, the Willamette River flows into the Columbia at Kelley Point Park, built on land acquired from the Port of Portland in 1984.
Lytton, British Columbia, Canada, is located at the confluence of the muddy Fraser River and the clearer Thompson River.
South America
Manaus, Brazil is on the Rio Negro near its confluence with the Amazon (see Meeting of Waters). It is the chief port and a hub for the region's extensive river system.
The Iguazú flows into the Paraná at the "Triple Frontier" (, ), the tripoint for Paraguay, Argentina, and Brazil.
In Ciudad Guayana, Venezuela there is a confluence between Orinoco River and Caroní River.
Confluences of non-rivers
Occasionally, "confluence" is used to describe the meeting of tidal or other non-riverine bodies of water, such as two canals or a canal and a lake. A one-mile (1.6 km) portion of the Industrial Canal in New Orleans accommodates the Gulf Intracoastal Waterway and the Mississippi River-Gulf Outlet Canal; therefore those three waterways are confluent there.
The term confluence can also apply to the process of merging or flowing together of other substance. For example, it may refer to the merger of the flow of two glaciers.
| Physical sciences | Hydrology | Earth science |
21281502 | https://en.wikipedia.org/wiki/Gun | Gun | A gun is a device designed to propel a projectile using pressure or explosive force. The projectiles are typically solid, but can also be pressurized liquid (e.g. in water guns/cannons), or gas (e.g. light-gas gun). Solid projectiles may be free-flying (as with bullets and artillery shells) or tethered (as with Tasers, spearguns and harpoon guns). A large-caliber gun is also called a cannon. Guns were designed as weapons for military use, and then found use in hunting. Now, there are guns, e.g., toy guns, water guns, paintball guns, etc., for many purposes.
The means of projectile propulsion vary according to designs, but are traditionally effected pneumatically by a high gas pressure contained within a barrel tube (gun barrel), produced either through the rapid exothermic combustion of propellants (as with firearms), or by mechanical compression (as with air guns). The high-pressure gas is introduced behind the projectile, pushing and accelerating it down the length of the tube, imparting sufficient launch velocity to sustain its further travel towards the target once the propelling gas ceases acting upon it after it exits the muzzle. Alternatively, new-concept linear motor weapons may employ an electromagnetic field to achieve acceleration, in which case the barrel may be substituted by guide rails (as in railguns) or wrapped with magnetic coils (as in coilguns).
The first devices identified as guns or proto-guns appeared in China from around AD 1000. By the end of the 13th century, they had become "true guns", metal barrel firearms that fired single projectiles which occluded the barrel. Gunpowder and gun technology spread throughout Eurasia during the 14th century.
Etymology and terminology
The origin of the English word gun is considered to derive from the name given to a particular historical weapon. Domina Gunilda was the name given to a remarkably large ballista, a mechanical bolt throwing weapon of enormous size, mounted at Windsor Castle during the 14th century. This name in turn may have derived from the Old Norse woman's proper name Gunnhildr which combines two Norse words referring to battle. "Gunnildr", which means "War-sword", was often shortened to "Gunna".
The earliest recorded use of the term "gonne" was in a Latin document . Other names for guns during this era were "schioppi" (Italian translation-"thunderers"), and "donrebusse" (Dutch translation-"thunder gun") which was incorporated into the English language as "blunderbuss". Artillerymen were often referred to as "gonners" and "artillers" "Hand gun" was first used in 1373 in reference to the handle of guns.
Definition
According to the Merriam-Webster dictionary, a gun could mean "a piece of ordnance usually with high muzzle velocity and comparatively flat trajectory," " a portable firearm," or "a device that throws a projectile."
Gunpowder and firearm historian Kenneth Chase defines "firearms" and "guns" in his Firearms: A Global History to 1700 as "gunpowder weapons that use the explosive force of the gunpowder to propel a projectile from a tube: cannons, muskets, and pistols are typical examples."
True gun
According to Tonio Andrade, a historian of gunpowder technology, a "true gun" is defined as a firearm which shoots a bullet that fits the barrel as opposed to one which does not, such as the shrapnel shooting fire lance. As such, the fire lance, which appeared between the 10th and 12th centuries AD, as well as other early metal barrel gunpowder weapons have been described as "proto-guns" Joseph Needham defined a type of firearm known as the "eruptor," which he described as a cross between a fire lance and a gun, as a "proto-gun" for the same reason. He defined a fully developed firearm, a "true gun," as possessing three basic features: a metal barrel, gunpowder with high nitrate content, and a projectile that occluded the barrel. The "true gun" appears to have emerged in late 1200s China, around 300 years after the appearance of the fire lance. Although the term "gun" postdates the invention of firearms, historians have applied it to the earliest firearms such as the Heilongjiang hand cannon of 1288 or the vase shaped European cannon of 1326.
Classic gun
Historians consider firearms to have reached the form of a "classic gun" in the 1480s, which persisted until the mid-18th century. This "classic" form displayed longer, lighter, more efficient, and more accurate design compared to its predecessors only 30 years prior. However this "classic" design changed very little for almost 300 years and cannons of the 1480s show little difference and surprising similarity with cannons later in the 1750s. This 300-year period during which the classic gun dominated gives it its moniker. The "classic gun" has also been described as the "modern ordnance synthesis."
History
Proto-gun
Gunpowder was invented in China during the 9th century. The first firearm was the fire lance, which was invented in China between the 10–12th centuries. It was depicted in a silk painting dated to the mid-10th century, but textual evidence of its use does not appear until 1132, describing the siege of De'an. It consisted of a bamboo tube of gunpowder tied to a spear or other polearm. By the late 1100s, ingredients such as pieces of shrapnel like porcelain shards or small iron pellets were added to the tube so that they would be blown out with the gunpowder. It was relatively short ranged and had a range of roughly 3 meters by the early 13th century. This fire lance is considered by some historians to be a "proto-gun" because its projectiles did not occlude the barrel. There was also another "proto-gun" called the eruptor, according to Joseph Needham, which did not have a lance but still did not shoot projectiles which occluded the barrel.
Transition to true guns
In due course, the proportion of saltpeter in the propellant was increased to maximise its explosive power. To better withstand that explosive power, the paper and bamboo of which fire-lance barrels were originally made came to be replaced by metal. And to take full advantage of that power, the shrapnel came to be replaced by projectiles whose size and shape filled the barrel more closely. Fire lance barrels made of metal appeared by 1276. Earlier in 1259 a pellet wad that filled the barrel was recorded to have been used as a fire lance projectile, making it the first recorded bullet in history. With this, the three basic features of a gun were put in place: a barrel made of metal, high-nitrate gunpowder, and a projectile which totally occludes the muzzle so that the powder charge exerts its full potential in propellant effect. The metal barrel fire lances began to be used without the lance and became guns by the late 13th century.
Guns such as the hand cannon were being used in the Yuan dynasty by the 1280s. Surviving cannons such as the Heilongjiang hand cannon and the Xanadu Gun have been found dating to the late 13th century and possibly earlier in the early 13th century.
In 1287, the Yuan dynasty deployed Jurchen troops with hand cannons to put down a rebellion by the Mongol prince Nayan. The History of Yuan records that the cannons of Li Ting's soldiers "caused great damage" and created "such confusion that the enemy soldiers attacked and killed each other." The hand cannons were used again in the beginning of 1288. Li Ting's "gun-soldiers" or chongzu () carried the hand cannons "on their backs". The passage on the 1288 battle is also the first to use the name chong () with the metal radical jin () for metal-barrel firearms. Chong was used instead of the earlier and more ambiguous term huo tong (fire tube; ), which may refer to the tubes of fire lances, proto-cannons, or signal flares. Hand cannons may have been used in the Mongol invasions of Japan. Japanese descriptions of the invasions mention iron and bamboo pao causing "light and fire" and emitting 2–3,000 iron bullets. The Nihon Kokujokushi, written around 1300, mentions huo tong (fire tubes) at the Battle of Tsushima in 1274 and the second coastal assault led by Holdon in 1281. The Hachiman Gudoukun of 1360 mentions iron pao "which caused a flash of light and a loud noise when fired." The Taiheki of 1370 mentions "iron pao shaped like a bell."
Spread
The exact nature of the spread of firearms and its route is uncertain. One theory is that gunpowder and cannons arrived in Europe via the Silk Road through the Middle East. Hasan al-Rammah had already written about fire lances in the 13th century, so proto-guns were known in the Middle East at that point. Another theory is that it was brought to Europe during the Mongol invasion in the first half of the 13th century.
The earliest depiction of a cannon in Europe dates to 1326 and evidence of firearm production can be found in the following year. The first recorded use of gunpowder weapons in Europe was in 1331 when two mounted German knights attacked Cividale del Friuli with gunpowder weapons of some sort. By 1338 hand cannons were in widespread use in France. English Privy Wardrobe accounts list "ribaldis", a type of cannon, in the 1340s, and siege guns were used by the English at Calais in 1346. Early guns and the men who used them were often associated with the devil and the gunner's craft was considered a black art, a point reinforced by the smell of sulfur on battlefields created from the firing of guns along with the muzzle blast and accompanying flash.
Around the late 14th century in Europe, smaller and portable hand-held cannons were developed, creating in effect the first smooth-bore personal firearm. In the late 15th century the Ottoman empire used firearms as part of its regular infantry. In the Middle East, the Arabs seem to have used the hand cannon to some degree during the 14th century. Cannons are attested in India starting from 1366.
The Joseon kingdom in Korea learned how to produce gunpowder from China by 1372 and started producing cannons by 1377. In Southeast Asia, Đại Việt soldiers used hand cannons at the very latest by 1390 when they employed them in killing Champa king Che Bong Nga. Chinese observer recorded the Javanese use of hand cannon for marriage ceremony in 1413 during Zheng He's voyage. Hand guns were utilized effectively during the Hussite Wars. Japan knew of gunpowder due to the Mongol invasions during the 13th century, but did not acquire a cannon until a monk took one back to Japan from China in 1510, and guns were not produced until 1543, when the Portuguese introduced matchlocks which were known as tanegashima to the Japanese.
Gunpowder technology entered Java in the Mongol invasion of Java (1293 A.D.). Majapahit under Mahapatih (prime minister) Gajah Mada utilized gunpowder technology obtained from the Yuan dynasty for use in the naval fleet. During the following years, the Majapahit army have begun producing cannons known as cetbang. Early cetbang (also called Eastern-style cetbang) resembled Chinese cannons and hand cannons. Eastern-style cetbangs were mostly made of bronze and were front-loaded cannons. It fires arrow-like projectiles, but round bullets and co-viative projectiles can also be used. These arrows can be solid-tipped without explosives, or with explosives and incendiary materials placed behind the tip. Near the rear, there is a combustion chamber or room, which refers to the bulging part near the rear of the gun, where the gunpowder is placed. The cetbang is mounted on a fixed mount, or as a hand cannon mounted on the end of a pole. There is a tube-like section on the back of the cannon. In the hand cannon-type cetbang, this tube is used as a socket for a pole.
Arquebus and musket
The arquebus was a firearm that appeared in Europe and the Ottoman Empire in the early 15th century. Its name is derived from the German word Hackenbüchse. It originally described a hand cannon with a lug or hook on the underside for stabilizing the weapon, usually on defensive fortifications. In the early 1500s, heavier variants known as "muskets" that were fired from resting Y-shaped supports appeared. The musket was able to penetrate heavy armor, and as a result armor declined, which also made the heavy musket obsolete. Although there is relatively little to no difference in design between arquebus and musket except in size and strength, it was the term musket which remained in use up into the 1800s. It may not be completely inaccurate to suggest that the musket was in its fabrication simply a larger arquebus. At least on one occasion the musket and arquebus have been used interchangeably to refer to the same weapon, and even referred to as an "arquebus musket." A Habsburg commander in the mid-1560s once referred to muskets as "double arquebuses."
A shoulder stock was added to the arquebus around 1470 and the matchlock mechanism sometime before 1475. The matchlock arquebus was the first firearm equipped with a trigger mechanism and the first portable shoulder-arms firearm. Before the matchlock, handheld firearms were fired from the chest, tucked under one arm, while the other arm maneuvered a hot pricker to the touch hole to ignite the gunpowder.
The Ottomans may have used arquebuses as early as the first half of the 15th century during the Ottoman–Hungarian wars of 1443–1444. The arquebus was used in substantial numbers during the reign of king Matthias Corvinus of Hungary (r. 1458–1490). Arquebuses were used by 1472 by the Spanish and Portuguese at Zamora. Likewise, the Castilians used arquebuses as well in 1476. Later, a larger arquebus known as a musket was used for breaching heavy armor, but this declined along with heavy armor. Matchlock firearms continued to be called musket. They were used throughout Asia by the mid-1500s.
Transition to classic guns
Guns reached their "classic" form in the 1480s. The "classic gun" is so called because of the long duration of its design, which was longer, lighter, more efficient, and more accurate compared to its predecessors 30 years prior. The design persisted for nearly 300 years and cannons of the 1480s show little variation from as well as surprising similarity with cannons three centuries later in the 1750s. This 300-year period during which the classic gun dominated gives it its moniker.
The classic gun differed from older generations of firearms through an assortment of improvements. Their longer length-to-bore ratio imparted more energy into the shot, enabling the projectile to shoot further. They were also lighter since the barrel walls were thinner, allowing faster dissipation of heat. They no longer needed the help of a wooden plug to load since they offered a tighter fit between projectile and barrel, further increasing the accuracy of firearms – and were deadlier due to developments such as gunpowder corning and iron shot.
Modern guns
Several developments in the 19th century led to the development of modern guns.
In 1815, Joshua Shaw invented percussion caps, which replaced the flintlock trigger system. The new percussion caps allowed guns to shoot reliably in any weather condition.
In 1835, Casimir Lefaucheux invented the first practical breech loading firearm with a cartridge. The new cartridge contained a conical bullet, a cardboard powder tube, and a copper base that incorporated a primer pellet.
Rifles
While rifled guns did exist prior to the 19th century in the form of grooves cut into the interior of a barrel, these were considered specialist weapons and limited in number.
The rate of fire of handheld guns began to increase drastically. In 1836, Johann Nicolaus von Dreyse invented the Dreyse needle gun, a breech-loading rifle which increased the rate of fire to six times that of muzzle loading weapons. In 1854, Volcanic Repeating Arms produced a rifle with a self-contained cartridge.
In 1849, Claude-Étienne Minié invented the Minié ball, the first projectile that could easily slide down a rifled barrel, which made rifles a viable military firearm, ending the smoothbore musket era. Rifles were deployed during the Crimean War with resounding success and proved vastly superior to smoothbore muskets.
In 1860, Benjamin Tyler Henry created the Henry rifle, the first reliable repeating rifle. An improved version of the Henry rifle was developed by Winchester Repeating Arms Company in 1873, known as the Model 1873 Winchester rifle.
Smokeless powder was invented in 1880 and began replacing gunpowder, which came to be known as black powder. By the start of the 20th century, smokeless powder was adopted throughout the world and black powder, what was previously known as gunpowder, was relegated to hobbyist usage.
Machine guns
In 1861, Richard Jordan Gatling invented the Gatling gun, the first successful machine gun, capable of firing 200 gunpowder cartridges in a minute. It was fielded by the Union forces during the American Civil War in the 1860s. In 1884, Hiram Maxim invented the Maxim gun, the first single-barreled machine gun.
The world's first submachine gun (a fully automatic firearm which fires pistol cartridges) able to be maneuvered by a single soldier is the MP 18.1, invented by Theodor Bergmann. It was introduced into service in 1918 by the German Army during World War I as the primary weapon of the Stosstruppen (assault groups specialized in trench combat).
In civilian use, the captive bolt pistol is used in agriculture to humanely stun farm animals for slaughter.
The first assault rifle was introduced during World War II by the Germans, known as the StG44. It was the first firearm to bridge the gap between long range rifles, machine guns, and short range submachine guns. Since the mid-20th century, guns that fire beams of energy rather than solid projectiles have been developed, and also guns that can be fired by means other than the use of gunpowder.
Operating principle
Most guns use compressed gas confined by the barrel to propel the bullet up to high speed, though devices operating in other ways are sometimes called guns. In firearms the high-pressure gas is generated by combustion, usually of gunpowder. This principle is similar to that of internal combustion engines, except that the bullet leaves the barrel, while the piston transfers its motion to other parts and returns down the cylinder. As in an internal combustion engine, the combustion propagates by deflagration rather than by detonation, and the optimal gunpowder, like the optimal motor fuel, is resistant to detonation. This is because much of the energy generated in detonation is in the form of a shock wave, which can propagate from the gas to the solid structure and heat or damage the structure, rather than staying as heat to propel the piston or bullet. The shock wave at such high temperature and pressure is much faster than that of any bullet, and would leave the gun as sound either through the barrel or the bullet itself rather than contributing to the bullet's velocity.
Components
Barrel
Barrel types include rifled—a series of spiraled grooves or angles within the barrel—when the projectile requires an induced spin to stabilize it, and smoothbore when the projectile is stabilized by other means or rifling is undesired or unnecessary. Typically, interior barrel diameter and the associated projectile size is a means to identify gun variations. Bore diameter is reported in several ways. The more conventional measure is reporting the interior diameter (bore) of the barrel in decimal fractions of the inch or in millimetres. Some guns—such as shotguns—report the weapon's gauge (which is the number of shot pellets having the same diameter as the bore produced from one English pound (454g) of lead) or—as in some British ordnance—the weight of the weapon's usual projectile.
Projectile
A gun projectile may be a simple, single-piece item like a bullet, a casing containing a payload like a shotshell or explosive shell, or complex projectile like a sub-caliber projectile and sabot. The propellant may be air, an explosive solid, or an explosive liquid. Some variations like the Gyrojet and certain other types combine the projectile and propellant into a single item.
Types
Military
Long gun
Arquebus
Blunderbuss
Musket
Musketoon
Wall gun
Grenade launcher
Submachine gun
Personal defense weapon
Rifle
Anti-materiel rifle
Anti-tank rifle
Automatic rifle
Lever-action rifle
Bolt-action rifle
Assault rifle
Battle rifle
Carbine
Designated marksman rifle
Service rifle
Semi-automatic rifle
Sniper rifle
Shotgun
Double-barreled shotgun
Sawed-off shotgun
Pump-action shotgun
Combat shotgun
Semi-automatic shotgun
Automatic shotgun
Handguns
Handgun
Derringer
Pistol
Semi-automatic pistol
Machine pistol
Service pistol
Revolver
Service revolver
Hunting
Air gun
BB gun
Elephant gun
Express rifle
Rimfire rifle
Speargun
Varmint rifle
Machine guns
Gatling gun
Rotary cannon
Rotary machine gun
Ribauldequin
Nordenfelt gun
Metal Storm
Volley gun
Mitrailleuse
Machine gun
General-purpose machine gun
Light machine gun
Squad Automatic Weapon
Infantry Automatic Rifle
Medium machine gun
Heavy machine gun
Autocannon
Autocannon
Chain gun
Revolver cannon
Artillery
Artillery gun
Cannon
Carronade
Falconet
Field gun
Howitzer
Mortar
Anti-aircraft gun
Anti-tank gun
Tank
Tank gun
Rescue equipment
Flare gun
Lyle gun
Training and entertainment
Airsoft gun
Cap gun
Drill Purpose Rifle
Nerf gun
Paintball gun
Potato cannon
Prop gun
Spud gun
Water gun
Directed-energy weapons
Anti-drone rifle
| Technology | Projectile weapons | null |
21281976 | https://en.wikipedia.org/wiki/Somatosensory%20system | Somatosensory system | The somatosensory system, or somatic sensory system is a subset of the sensory nervous system. It has two subdivisions, one for the detection of mechanosensory information related to touch, and the other for the nociception detection of pain and temperature. The main functions of the somatosensory system are the perception of external stimuli, the perception of internal stimuli, and the regulation of body position and balance (proprioception).
Mechanosensory information includes that of light touch, vibration, pressure and tension in the skin. Much of this information belongs to the sense of touch which is a general somatic sense in contrast to the special senses of sight, smell, taste, hearing, and balance.
Nociceptory information is that received from pain and temperature that is deemed as harmful (noxious). Thermoreceptors relay temperature information in normal circumstances. Nociceptors are specialised receptors for signals of pain.
The sense of touch in perceiving the environment uses special sensory receptors in the skin called cutaneous receptors. They include mechanoreceptors such as tactile corpuscles that relay information about pressure and vibration; nociceptors, and thermoreceptors for temperature perception.
Stimulation of the receptors activate peripheral sensory neurons that convey signals to the spinal cord that may drive a responsive reflex, and may also be conveyed to the brain for conscious perception. Somatosensory information from the face and head enter the brain via cranial nerves such as the trigeminal nerve.
The neural pathways that go to the brain are structured such that information about the location of the physical stimulus is preserved. In this way, neighboring neurons in the somatosensory cortex represent nearby locations on the skin or in the body, creating a map or sensory homunculus.
Touch communication
Tactile signing
Tactile signing is a common means of communication used by people with deafblindness. It is based on a sign language or another system of manual communication.
Emotion communication
Humans can communicate specific emotions through touch alone including anger, fear, disgust, love, gratitude, and sympathy via touch at much-better-than-chance levels.
Overview
Sensory receptors
The two different types of mechanoreceptor in the skin are termed low-threshold mechanoreceptors, and high threshold mechanoreceptors.
The four mechanoreceptors in glabrous skin are low-threshold that respond to harmless stimuli. They are innervated by four different afferent fibers. High-threshold mechanoreceptors, respond to harmful stimuli.
Merkel cell nerve endings are found in the basal epidermis and hair follicles; they react to low vibrations (5–15 Hz) and deep static touch such as shapes and edges. Due to having a small receptive field (extremely detailed information), they are used in areas like fingertips the most; they are not covered (shelled) and thus respond to pressures over long periods.
Tactile corpuscles react to moderate vibration (10–50 Hz) and light touch. They are located in the dermal papillae; due to their reactivity, they are primarily located in fingertips and lips. They respond in quick action potentials, unlike Merkel nerve endings. They are responsible for the ability to read Braille and feel gentle stimuli.
Pacinian corpuscles determine gross touch and distinguish rough and soft substances. They react in quick action potentials, especially to vibrations around 250 Hz (even up to centimeters away). They are the most sensitive to vibrations and have large receptor fields. Pacinian corpuscles react only to sudden stimuli so pressures like clothes that are always compressing their shape are quickly ignored. They have also been implicated in detecting the location of touch sensations on handheld tools.
Bulbous corpuscles react slowly and respond to sustained skin stretch. They are responsible for the feeling of object slippage and play a major role in the kinesthetic sense and control of finger position and movement. Merkel and bulbous cells - slow-response - are myelinated; the rest - fast-response - are not. All of these receptors are activated upon pressures that distort their shape causing an action potential.
Somatosensory cortex
The postcentral gyrus is in the parietal lobe and its cortex is the primary somatosensory cortex (Brodmann areas 3, 2 and 1) collectively referred to as S1.
BA3 receives the densest projections from the thalamus. BA3a is involved with the sense of relative position of neighboring body parts and amount of effort being used during movement. BA3b is responsible for distributing somatosensory information, it projects texture information to BA1 and shape and size information to BA2.
Region S2 (secondary somatosensory cortex) divides into Area S2 and parietal ventral area. Area S2 is involved with specific touch perception and is thus integrally linked with the amygdala and hippocampus to encode and reinforce memories.
Parietal ventral area is the somatosensory relay to the premotor cortex and somatosensory memory hub, BA5.
BA5 is the topographically organized somato memory field and association area.
BA1 processes texture info while BA2 processes size and shape information.
Area S2 processes light touch, pain, visceral sensation, and tactile attention.
S1 processes the remaining info (crude touch, pain, temperature).
BA7 integrates visual and proprioceptive info to locate objects in space.
The insular cortex (insula) plays a role in the sense of bodily-ownership, bodily self-awareness, and perception. Insula also plays a role in conveying info about sensual touch, pain, temperature, itch, and local oxygen status. Insula is a highly connected relay and thus is involved in numerous functions.
Structure
The somatosensory system is spread through all major parts of the vertebrate body. It consists both of sensory receptors and sensory neurons in the periphery (skin, muscle and organs for example), to deeper neurons within the central nervous system.
General somatosensory pathway
All afferent touch/vibration information ascends the spinal cord via the dorsal column-medial lemniscus pathway via gracilis (T7 and below) or cuneatus (T6 and above). Cuneatus sends signals to the cochlear nucleus indirectly via spinal grey matter, this info is used in determining if a perceived sound is just villi noise/irritation. All fibers cross (left becomes right) in the medulla.
A somatosensory pathway will typically have three neurons: first-order, second-order, and third-order.
The first-order neuron is a type of pseudounipolar neuron and always has its cell body in the dorsal root ganglion of the spinal nerve with a peripheral axon innervating touch mechanoreceptors and a central axon synapsing on the second-order neuron. If the somatosensory pathway is in parts of the head or neck not covered by the cervical nerves, the first-order neuron will be the trigeminal nerve ganglia or the ganglia of other sensory cranial nerves).
The second-order neuron has its cell body either in the spinal cord or in the brainstem. This neuron's ascending axons will cross (decussate) to the opposite side either in the spinal cord or in the brainstem.
In the case of touch and certain types of pain, the third-order neuron has its cell body in the ventral posterior nucleus of the thalamus and ends in the postcentral gyrus of the parietal lobe in the primary somatosensory cortex (or S1).
Photoreceptors, similar to those found in the retina of the eye, detect potentially damaging ultraviolet radiation (ultraviolet A specifically), inducing increased production of melanin by melanocytes. Thus tanning potentially offers the skin rapid protection from DNA damage and sunburn caused by ultraviolet radiation (DNA damage caused by ultraviolet B). However, whether this offers protection is debatable, because the amount of melanin released by this process is modest in comparison to the amounts released in response to DNA damage caused by ultraviolet B radiation.
Tactile feedback
The tactile feedback from proprioception is derived from the proprioceptors in the skin, muscles, and joints.
Balance
The receptor for the sense of balance resides in the vestibular system in the ear (for the three-dimensional orientation of the head, and by inference, the rest of the body). Balance is also mediated by the kinesthetic reflex fed by proprioception (which senses the relative location of the rest of the body to the head). In addition, proprioception estimates the location of objects which are sensed by the visual system (which provides confirmation of the place of those objects relative to the body), as input to the mechanical reflexes of the body.
Fine touch and crude touch
Fine touch (or discriminative touch) is a sensory modality that allows a subject to sense and localize touch. The form of touch where localization is not possible is known as crude touch. The dorsal column–medial lemniscus pathway is the pathway responsible for the sending of fine touch information to the cerebral cortex of the brain.
Crude touch (non-discriminating) is a sensory modality that allows the subject to sense that something has touched them, without being able to localize where they were touched (contrasting "fine touch"). Its fibres are carried in the spinothalamic tract, unlike the fine touch, which is carried in the dorsal column.
As fine touch normally works in parallel to crude touch, a person will be able to localize touch until fibres carrying fine touch (in the dorsal column–medial lemniscus pathway) have been disrupted. Then the subject will feel the touch, but be unable to identify where they were touched.
Neural processing of social touch
The somatosensory cortex encodes incoming sensory information from receptors all over the body. Affective touch is a type of sensory information that elicits an emotional reaction and is usually social in nature, such as a physical human touch. This type of information is actually coded differently than other sensory information. Intensity of affective touch is still encoded in the primary somatosensory cortex and is processed in a similar way to emotions invoked by sight and sound, as exemplified by the increase of adrenaline caused by the social touch of a loved one, as opposed to the physical inability to touch someone you do not love.
Meanwhile, the feeling of pleasantness associated with affective touch activates the anterior cingulate cortex more than the primary somatosensory cortex. Functional magnetic resonance imaging (fMRI) data shows that increased blood-oxygen-level contrast (BOLD) signal in the anterior cingulate cortex as well as the prefrontal cortex is highly correlated with pleasantness scores of an affective touch. Inhibitory transcranial magnetic stimulation (TMS) of the primary somatosensory cortex inhibits the perception of affective touch intensity, but not affective touch pleasantness. Therefore, the S1 is not directly involved in processing socially affective touch pleasantness, but still plays a role in discriminating touch location and intensity.
Tactile interaction is important amongst some animals. Usually, tactile contact between two animals occurs through stroking, licking, or grooming. These behaviours are essential for the individual's social healthcare, as in the hypothalamus they induce the release of oxytocin, a hormone that decreases stress and anxiety and increases social bonding between animals.
More precisely, the consistency of oxytocin neuron activation in rats stroked by humans has been observed, especially in the caudal paraventricular nucleus. It was found that this affiliative relationship induced by tactile contact is common no matter the relationship between the two individuals (mother-infant, male-female, human-animal). It has also been discovered that the level of oxytocin release through this behaviour correlates with the time course of social interaction as longer stroking induced a greater release of the hormone.
The importance of somatosensory stimulation in social animals such as primates has also been observed. Grooming is part of the social interaction primates exert on their conspecifics. This interaction is required between individuals to maintain the affiliative relationship within the group, avoid internal conflict and increase group bonding. However, such social interaction requires the recognition of every member in the group. As such, it has been observed that the size of the neocortex is positively correlated with the size of the group, reflecting a limit to the number of recognizable members amongst which grooming can occur. Furthermore, the time course of grooming is related to vulnerability due to predation to which animals are exposed to whilst performing such social interaction. The relationship between tactile interaction, stress reduction and social bonding depends on the evaluation of risks that occur during conducting such behaviours in the wild life, and further research is required to unveil the connection between tactile caring and fitness level.
Studies show a correlation between touching a soft or hard object and how a person thinks or even makes decisions. Further, between the firmness of a touch and the evoking of gender stereotyping.
Tactile memories as part of haptic memory, are organized somatotopically, following the organization of the somatosensory cortex.
Individual variation
A variety of studies have measured and investigated the causes for differences between individuals in the sense of fine touch. One well-studied area is passive tactile spatial acuity, the ability to resolve the fine spatial details of an object pressed against the stationary skin. A variety of methods have been used to measure passive tactile spatial acuity, perhaps the most rigorous being the grating orientation task. In this task subjects identify the orientation of a grooved surface presented in two different orientations, which can be applied manually or with automated equipment. Many studies have shown a decline in passive tactile spatial acuity with age; the reasons for this decline are unknown, but may include loss of tactile receptors during normal aging. Remarkably, index finger passive tactile spatial acuity is better among adults with smaller index fingertips; this effect of finger size has been shown to underlie the better passive tactile spatial acuity of women, on average, compared to men. The density of tactile corpuscles, a type of mechanoreceptor that detects low-frequency vibrations, is greater in smaller fingers; the same may hold for Merkel cells, which detect the static indentations important for fine spatial acuity. Among children of the same age, those with smaller fingers also tend to have better tactile acuity. Many studies have shown that passive tactile spatial acuity is enhanced among blind individuals compared to sighted individuals of the same age, possibly because of cross modal plasticity in the cerebral cortex of blind individuals. Perhaps also due to cortical plasticity, individuals who have been blind since birth reportedly consolidate tactile information more rapidly than sighted people.
Clinical significance
A somatosensory deficiency may be caused by a peripheral neuropathy involving peripheral nerves of the somatosensory system. This may present as numbness or paresthesia.
Society and culture
Haptic technology can provide touch sensation in virtual and real environments. In the field of speech therapy, tactile feedback can be used to treat speech disorders.
Affectionate touch is present in everyday life and can take multiple forms. These actions, however, seem to carry specific functions even though the evolutionary benefit from such a wide range of behaviours is not entirely understood. Researchers investigated the expression patterns and characteristics of 8 different affectionate touch actions - embracing, holding, kissing, leaning, petting, squeezing, stroking, and tickling - in a self-report study. It was found that the affectionate touch has distinct target areas on the body, different associated affect, comfort-value, and expression frequency based on the type of touch action that is performed.
Besides the rather obvious sensory consequences of touch, it can also affect higher-level aspects of cognition such as social judgements and decision-making. This effect might arise due to a physical-to-mental scaffolding process in early development, whereby sensorimotor experiences are linked to the emergence of conceptual knowledge. Such links might be maintained throughout life, and so touching an object may cue the physical sensation to its related conceptual processing. Indeed, it was found that different physical properties - weight, texture, and hardness - of a touched object can influence social judgement and decision-making. For example, participants described a passage of a social interaction to be harsher when they touched a hard wooden block instead of a soft blanket prior to the task. Building on these findings, the ability of touch to have an unconscious influence on such higher-order thoughts may provide a novel tool for marketing and communication strategies.
| Biology and health sciences | Nervous system | null |
21282020 | https://en.wikipedia.org/wiki/Hearing | Hearing | Hearing, or auditory perception, is the ability to perceive sounds through an organ, such as an ear, by detecting vibrations as periodic changes in the pressure of a surrounding medium. The academic field concerned with hearing is auditory science.
Sound may be heard through solid, liquid, or gaseous matter. It is one of the traditional five senses. Partial or total inability to hear is called hearing loss.
In humans and other vertebrates, hearing is performed primarily by the auditory system: mechanical waves, known as vibrations, are detected by the ear and transduced into nerve impulses that are perceived by the brain (primarily in the temporal lobe). Like touch, audition requires sensitivity to the movement of molecules in the world outside the organism. Both hearing and touch are types of mechanosensation.
Hearing mechanism
There are three main components of the human auditory system: the outer ear, the middle ear, and the inner ear.
Outer ear
The outer ear includes the pinna, the visible part of the ear, as well as the ear canal, which terminates at the eardrum, also called the tympanic membrane. The pinna serves to focus sound waves through the ear canal toward the eardrum. Because of the asymmetrical character of the outer ear of most mammals, sound is filtered differently on its way into the ear depending on the location of its origin. This gives these animals the ability to localize sound vertically. The eardrum is an airtight membrane, and when sound waves arrive there, they cause it to vibrate following the waveform of the sound. Cerumen (ear wax) is produced by ceruminous and sebaceous glands in the skin of the human ear canal, protecting the ear canal and tympanic membrane from physical damage and microbial invasion.
Middle ear
The middle ear consists of a small air-filled chamber that is located medial to the eardrum. Within this chamber are the three smallest bones in the body, known collectively as the ossicles which include the malleus, incus, and stapes (also known as the hammer, anvil, and stirrup, respectively). They aid in the transmission of the vibrations from the eardrum into the inner ear, the cochlea. The purpose of the middle ear ossicles is to overcome the impedance mismatch between air waves and cochlear waves, by providing impedance matching.
Also located in the middle ear are the stapedius muscle and tensor tympani muscle, which protect the hearing mechanism through a stiffening reflex. The stapes transmits sound waves to the inner ear through the oval window, a flexible membrane separating the air-filled middle ear from the fluid-filled inner ear. The round window, another flexible membrane, allows for the smooth displacement of the inner ear fluid caused by the entering sound waves.
Inner ear
The inner ear consists of the cochlea, which is a spiral-shaped, fluid-filled tube. It is divided lengthwise by the organ of Corti, which is the main organ of mechanical to neural transduction. Inside the organ of Corti is the basilar membrane, a structure that vibrates when waves from the middle ear propagate through the cochlear fluid – endolymph. The basilar membrane is tonotopic, so that each frequency has a characteristic place of resonance along it. Characteristic frequencies are high at the basal entrance to the cochlea, and low at the apex. Basilar membrane motion causes depolarization of the hair cells, specialized auditory receptors located within the organ of Corti. While the hair cells do not produce action potentials themselves, they release neurotransmitter at synapses with the fibers of the auditory nerve, which does produce action potentials. In this way, the patterns of oscillations on the basilar membrane are converted to spatiotemporal patterns of firings which transmit information about the sound to the brainstem.
Neuronal
The sound information from the cochlea travels via the auditory nerve to the cochlear nucleus in the brainstem. From there, the signals are projected to the inferior colliculus in the midbrain tectum. The inferior colliculus integrates auditory input with limited input from other parts of the brain and is involved in subconscious reflexes such as the auditory startle response.
The inferior colliculus in turn projects to the medial geniculate nucleus, a part of the thalamus where sound information is relayed to the primary auditory cortex in the temporal lobe. Sound is believed to first become consciously experienced at the primary auditory cortex. Around the primary auditory cortex lies Wernickes area, a cortical area involved in interpreting sounds that is necessary to understand spoken words.
Disturbances (such as stroke or trauma) at any of these levels can cause hearing problems, especially if the disturbance is bilateral. In some instances it can also lead to auditory hallucinations or more complex difficulties in perceiving sound.
Hearing tests
Hearing can be measured by behavioral tests using an audiometer. Electrophysiological tests of hearing can provide accurate measurements of hearing thresholds even in unconscious subjects. Such tests include auditory brainstem evoked potentials (ABR), otoacoustic emissions (OAE) and electrocochleography (ECochG). Technical advances in these tests have allowed hearing screening for infants to become widespread.
Hearing can be measured by mobile applications which includes audiological hearing test function or hearing aid application. These applications allow the user to measure hearing thresholds at different frequencies (audiogram). Despite possible errors in measurements, hearing loss can be detected.
Hearing loss
There are several different types of hearing loss: conductive hearing loss, sensorineural hearing loss and mixed types. Recently, the term of Aural Diversity has come into greater use, to communicate hearing loss and differences in a less negatively-associated term.
There are defined degrees of hearing loss:
Mild hearing loss - People with mild hearing loss have difficulties keeping up with conversations, especially in noisy surroundings. The most quiet sounds that people with mild hearing loss can hear with their better ear are between 25 and 40 dB HL.
Moderate hearing loss - People with moderate hearing loss have difficulty keeping up with conversations when they are not using a hearing aid. On average, the most quiet sounds heard by people with moderate hearing loss with their better ear are between 40 and 70 dB HL.
Severe hearing loss - People with severe hearing loss depend on powerful hearing aid. However, they often rely on lip-reading even when they are using hearing aids. The most quiet sounds heard by people with severe hearing loss with their better ear are between 70 and 95 dB HL.
Profound hearing loss - People with profound hearing loss are very hard of hearing and they mostly rely on lip-reading and sign language. The most quiet sounds heard by people with profound hearing loss with their better ear are from 95 dB HL or more.
Causes
Heredity
Congenital conditions
Presbycusis
Acquired
Noise-induced hearing loss
Ototoxic drugs and chemicals
Infection
Prevention
Hearing protection is the use of devices designed to prevent noise-induced hearing loss (NIHL), a type of post-lingual hearing impairment. The various means used to prevent hearing loss generally focus on reducing the levels of noise to which people are exposed. One way this is done is through environmental modifications such as acoustic quieting, which may be achieved with as basic a measure as lining a room with curtains, or as complex a measure as employing an anechoic chamber, which absorbs nearly all sound. Another means is the use of devices such as earplugs, which are inserted into the ear canal to block noise, or earmuffs, objects designed to cover a person's ears entirely.
Management
The loss of hearing, when it is caused by neural loss, cannot presently be cured. Instead, its effects can be mitigated by the use of audioprosthetic devices, i.e. hearing assistive devices such as hearing aids and cochlear implants. In a clinical setting, this management is offered by otologists and audiologists.
Relation to health
Hearing loss is associated with Alzheimer's disease and dementia with a greater degree of hearing loss tied to a higher risk. There is also an association between type 2 diabetes and hearing loss.
Hearing underwater
Hearing threshold and the ability to localize sound sources are reduced underwater in humans, but not in aquatic animals, including whales, seals, and fish which have ears adapted to process water-borne sound.
In vertebrates
Not all sounds are normally audible to all animals. Each species has a range of normal hearing for both amplitude and frequency. Many animals use sound to communicate with each other, and hearing in these species is particularly important for survival and reproduction. In species that use sound as a primary means of communication, hearing is typically most acute for the range of pitches produced in calls and speech.
Frequency range
Frequencies capable of being heard by humans are called audio or sonic. The range is typically considered to be between 20 Hz and 20,000 Hz. Frequencies higher than audio are referred to as ultrasonic, while frequencies below audio are referred to as infrasonic. Some bats use ultrasound for echolocation while in flight. Dogs are able to hear ultrasound, which is the principle of 'silent' dog whistles. Snakes sense infrasound through their jaws, and baleen whales, giraffes, dolphins and elephants use it for communication. Some fish have the ability to hear more sensitively due to a well-developed, bony connection between the ear and their swim bladder. This "aid to the deaf" for fishes appears in some species such as carp and herring.
Time discrimination
Human perception of audio signal time separation has been measured to less than 10 microseconds (10μs). This does not mean that frequencies above are audible, but that time discrimination is not directly coupled with frequency range. Georg Von Békésy in 1929 identifying sound source directions suggested humans can resolve timing differences of 10μs or less. In 1976 Jan Nordmark's research indicated inter-aural resolution better than 2μs. Milind Kuncher's 2007 research resolved time misalignment to under 10μs.
In birds
In invertebrates
Even though they do not have ears, invertebrates have developed other structures and systems to decode vibrations traveling through the air, or “sound”. Charles Henry Turner was the first scientist to formally show this phenomenon through rigorously controlled experiments in ants. Turner ruled out the detection of ground vibration and suggested that other insects likely have auditory systems as well.
Many insects detect sound through the way air vibrations deflect hairs along their body. Some insects have even developed specialized hairs tuned to detecting particular frequencies, such as certain caterpillar species that have evolved hair with properties such that it resonates most with the sound of buzzing wasps, thus warning them of the presence of natural enemies.
Some insects possess a tympanal organ. These are "eardrums", that cover air filled chambers on the legs. Similar to the hearing process with vertebrates, the eardrums react to sonar waves. Receptors that are placed on the inside translate the oscillation into electric signals and send them to the brain. Several groups of flying insects that are preyed upon by echolocating bats can perceive the ultrasound emissions this way and reflexively practice ultrasound avoidance.
| Biology and health sciences | Nervous system | null |
21282070 | https://en.wikipedia.org/wiki/Taste | Taste | The gustatory system or sense of taste is the sensory system that is partially responsible for the perception of taste. Taste is the perception stimulated when a substance in the mouth reacts chemically with taste receptor cells located on taste buds in the oral cavity, mostly on the tongue. Taste, along with the sense of smell and trigeminal nerve stimulation (registering texture, pain, and temperature), determines flavors of food and other substances. Humans have taste receptors on taste buds and other areas, including the upper surface of the tongue and the epiglottis. The gustatory cortex is responsible for the perception of taste.
The tongue is covered with thousands of small bumps called papillae, which are visible to the naked eye. Within each papilla are hundreds of taste buds. The exceptions to this is the filiform papillae that do not contain taste buds. There are between 2000 and 5000 taste buds that are located on the back and front of the tongue. Others are located on the roof, sides and back of the mouth, and in the throat. Each taste bud contains 50 to 100 taste receptor cells.
Taste receptors in the mouth sense the five basic tastes: sweetness, sourness, saltiness, bitterness, and savoriness (also known as savory or umami). Scientific experiments have demonstrated that these five tastes exist and are distinct from one another. Taste buds are able to tell different tastes apart when they interact with different molecules or ions. Sweetness, savoriness, and bitter tastes are triggered by the binding of molecules to G protein-coupled receptors on the cell membranes of taste buds. Saltiness and sourness are perceived when alkali metals or hydrogen ions meet taste buds, respectively.
The basic tastes contribute only partially to the sensation and flavor of food in the mouth—other factors include smell, detected by the olfactory epithelium of the nose; texture, detected through a variety of mechanoreceptors, muscle nerves, etc.; temperature, detected by temperature receptors; and "coolness" (such as of menthol) and "hotness" (pungency), by chemesthesis.
As the gustatory system senses both harmful and beneficial things, all basic tastes bring either caution or craving depending upon the effect the things they sense have on the body. Sweetness helps to identify energy-rich foods, while bitterness warns people of poisons.
Among humans, taste perception begins to fade during ageing, tongue papillae are lost, and saliva production slowly decreases. Humans can also have distortion of tastes (dysgeusia). Not all mammals share the same tastes: some rodents can taste starch (which humans cannot), cats cannot taste sweetness, and several other carnivores, including hyenas, dolphins, and sea lions, have lost the ability to sense up to four of their ancestral five basic tastes.
Basic tastes
The gustatory system allows animals to distinguish between safe and harmful food and to gauge different foods' nutritional value. Digestive enzymes in saliva begin to dissolve food into base chemicals that are washed over the papillae and detected as tastes by the taste buds. The tongue is covered with thousands of small bumps called papillae, which are visible to the naked eye. Within each papilla are hundreds of taste buds. The exception to this are the filiform papillae, which do not contain taste buds. There are between 2,000 and 5,000 taste buds that are located on the back and front of the tongue. Others are located on the roof, sides and back of the mouth, and in the throat. Each taste bud contains 50 to 100 taste-receptor cells.
The five specific tastes received by taste receptors are saltiness, sweetness, bitterness, sourness, and savoriness (often known by its Japanese name , which translates to 'deliciousness').
As of the early 20th century, Western physiologists and psychologists believed that there were four basic tastes: sweetness, sourness, saltiness, and bitterness. The concept of a "savory" taste was not present in Western science at that time, but was postulated in Japanese research.
One study found that salt and sour taste mechanisms both detect, in different ways, the presence of sodium chloride (salt) in the mouth. Acids are also detected and perceived as sour. The detection of salt is important to many organisms, but especially mammals, as it serves a critical role in ion and water homeostasis in the body. It is specifically needed in the mammalian kidney as an osmotically active compound that facilitates passive re-uptake of water into the blood. Because of this, salt elicits a pleasant taste in most humans.
Sour and salt tastes can be pleasant in small quantities, but in larger quantities become more and more unpleasant to taste. For sour taste, this presumably is because the sour taste can signal under-ripe fruit, rotten meat, and other spoiled foods, which can be dangerous to the body because of bacteria that grow in such media. Additionally, sour taste signals acids, which can cause serious tissue damage.
Sweet taste signals the presence of carbohydrates in solution. Since carbohydrates have a very high calorie count (saccharides have many bonds, therefore much energy), they are desirable to the human body, which evolved to seek out the highest-calorie-intake foods. They are used as direct energy (sugars) and storage of energy (glycogen). Many non-carbohydrate molecules trigger a sweet response, leading to the development of many artificial sweeteners, including saccharin, sucralose, and aspartame. It is still unclear how these substances activate the sweet receptors and what adaptative significance this has had.
The savory taste (known in Japanese as ), identified by Japanese chemist Kikunae Ikeda, signals the presence of the amino acid L-glutamate. The amino acids in proteins are used in the body to build muscles and organs, and to transport molecules (hemoglobin), antibodies, and the organic catalysts known as enzymes. These are all critical molecules, and it is important to have a steady supply of amino acids; consequently, savory tastes trigger a pleasurable response, encouraging the intake of peptides and proteins.
Pungency (piquancy or hotness) had traditionally been considered a sixth basic taste. In 2015, researchers suggested a new basic taste of fatty acids called "fat taste", although "oleogustus" and "pinguis" have both been proposed as alternate terms.
Sweetness
Sweetness, usually regarded as a pleasurable sensation, is produced by the presence of sugars and substances that mimic sugar. Sweetness may be connected to aldehydes and ketones, which contain a carbonyl group. Sweetness is detected by a variety of G protein coupled receptors (GPCR) coupled to the G protein gustducin found on the taste buds. At least two different variants of the "sweetness receptors" must be activated for the brain to register sweetness. Compounds the brain senses as sweet are compounds that can bind with varying bond strength to two different sweetness receptors. These receptors are T1R2+3 (heterodimer) and T1R3 (homodimer), which account for all sweet sensing in humans and animals.
Taste detection thresholds for sweet substances are rated relative to sucrose, which has an index of 1. The average human detection threshold for sucrose is 10 millimoles per liter. For lactose it is 30 millimoles per liter, with a sweetness index of 0.3, and 5-nitro-2-propoxyaniline 0.002 millimoles per liter. "Natural" sweeteners such as saccharides activate the GPCR, which releases gustducin. The gustducin then activates the molecule adenylate cyclase, which catalyzes the production of the molecule cAMP, or adenosine 3', 5'-cyclic monophosphate. This molecule closes potassium ion channels, leading to depolarization and neurotransmitter release. Synthetic sweeteners such as saccharin activate different GPCRs and induce taste receptor cell depolarization by an alternate pathway.
Sourness
Sourness is the taste that detects acidity. The sourness of substances is rated relative to dilute hydrochloric acid, which has a sourness index of 1. By comparison, tartaric acid has a sourness index of 0.7, citric acid an index of 0.46, and carbonic acid an index of 0.06.
Sour taste is detected by a small subset of cells that are distributed across all taste buds called Type III taste receptor cells. H+ ions (protons) that are abundant in sour substances can directly enter the Type III taste cells through a proton channel. This channel was identified in 2018 as otopetrin 1 (OTOP1). The transfer of positive charge into the cell can itself trigger an electrical response. Some weak acids such as acetic acid can also penetrate taste cells; intracellular hydrogen ions inhibit potassium channels, which normally function to hyperpolarize the cell. By a combination of direct intake of hydrogen ions through OTOP1 ion channels (which itself depolarizes the cell) and the inhibition of the hyperpolarizing channel, sourness causes the taste cell to fire action potentials and release neurotransmitter.
The most common foods with natural sourness are fruits, such as lemon, lime, grape, orange, tamarind, and bitter melon. Fermented foods, such as wine, vinegar or yogurt, may have sour taste. Children show a greater enjoyment of sour flavors than adults, and sour candy containing citric acid or malic acid is common.
Saltiness
Saltiness taste seems to have two components: a low-salt signal and a high-salt signal. The low-salt signal causes a sensation of deliciousness, while the high-salt signal typically causes the sensation of "too salty".
The low-salt signal is understood to be caused by the epithelial sodium channel (ENaC), which is composed of three subunits. ENaC in the taste cells allow sodium cations to enter the cell. This on its own depolarizes the cell, and opens voltage-dependent calcium channels, flooding the cell with positive calcium ions and leading to neurotransmitter release. ENaC can be blocked by the drug amiloride in many mammals, especially rats. The sensitivity of the low-salt taste to amiloride in humans is much less pronounced, leading to conjecture that there may be additional low-salt receptors besides ENaC to be discovered.
A number of similar cations also trigger the low salt signal. The size of lithium and potassium ions most closely resemble those of sodium, and thus the saltiness is most similar. In contrast, rubidium and caesium ions are far larger, so their salty taste differs accordingly. The saltiness of substances is rated relative to sodium chloride (NaCl), which has an index of 1. Potassium, as potassium chloride (KCl), is the principal ingredient in salt substitutes and has a saltiness index of 0.6.
Other monovalent cations, e.g. ammonium (NH4+), and divalent cations of the alkali earth metal group of the periodic table, e.g. calcium (Ca2+), ions generally elicit a bitter rather than a salty taste even though they, too, can pass directly through ion channels in the tongue, generating an action potential. But the chloride of calcium is saltier and less bitter than potassium chloride, and is commonly used in pickle brine instead of KCl.
The high-salt signal is still very poorly understood as of 2023. Even in rodents, this signal is not blocked by amiloride. Sour and bitter cells trigger on high chloride levels, but the specific receptor is still being identified.
Bitterness
Bitterness is one of the most sensitive of the tastes, and many perceive it as unpleasant, sharp, or disagreeable, but it is sometimes desirable and intentionally added via various bittering agents. Common bitter foods and beverages include coffee, unsweetened cocoa, South American mate, coca tea, bitter gourd, uncured olives, citrus peel, some varieties of cheese, many plants in the family Brassicaceae, dandelion greens, horehound, wild chicory, and escarole. The ethanol in alcoholic beverages tastes bitter, as do the additional bitter ingredients found in some alcoholic beverages including hops in beer and gentian in bitters. Quinine is also known for its bitter taste and is found in tonic water.
Bitterness is of interest to those who study evolution, as well as various health researchers since a large number of natural bitter compounds are known to be toxic. The ability to detect bitter-tasting, toxic compounds at low thresholds is considered to provide an important protective function. Plant leaves often contain toxic compounds, and among leaf-eating primates there is a tendency to prefer immature leaves, which tend to be higher in protein and lower in fiber and poisons than mature leaves. Amongst humans, various food processing techniques are used worldwide to detoxify otherwise inedible foods and make them palatable. Furthermore, the use of fire, changes in diet, and avoidance of toxins has led to neutral evolution in human bitter sensitivity. This has allowed several loss of function mutations that has led to a reduced sensory capacity towards bitterness in humans when compared to other species.
The threshold for stimulation of bitter taste by quinine averages a concentration of 8 μM (8 micromolar). The taste thresholds of other bitter substances are rated relative to quinine, which is thus given a reference index of 1. For example, brucine has an index of 11, is thus perceived as intensely more bitter than quinine, and is detected at a much lower solution threshold. The most bitter natural substance is amarogentin, a compound present in the roots of the plant Gentiana lutea, and the most bitter substance known is the synthetic chemical denatonium, which has an index of 1,000. It is used as an aversive agent (a bitterant) that is added to toxic substances to prevent accidental ingestion. It was discovered accidentally in 1958 during research on a local anesthetic by T. & H. Smith of Edinburgh, Scotland.
Research has shown that TAS2Rs (taste receptors, type 2, also known as T2Rs) such as TAS2R38 coupled to the G protein gustducin are responsible for the human ability to taste bitter substances. They are identified not only by their ability to taste for certain "bitter" ligands, but also by the morphology of the receptor itself (surface bound, monomeric). The TAS2R family in humans is thought to comprise about 25 different taste receptors, some of which can recognize a wide variety of bitter-tasting compounds. Over 670 bitter-tasting compounds have been identified, on a bitter database, of which over 200 have been assigned to one or more specific receptors. It is speculated that the selective constraints on the TAS2R family have been weakened due to the relatively high rate of mutation and pseudogenization. Researchers use two synthetic substances, phenylthiocarbamide (PTC) and 6-n-propylthiouracil (PROP) to study the genetics of bitter perception. These two substances taste bitter to some people, but are virtually tasteless to others. Among the tasters, some are so-called "supertasters" to whom PTC and PROP are extremely bitter. The variation in sensitivity is determined by two common alleles at the TAS2R38 locus. This genetic variation in the ability to taste a substance has been a source of great interest to those who study genetics.
Gustducin is made of three subunits. When it is activated by the GPCR, its subunits break apart and activate phosphodiesterase, a nearby enzyme, which in turn converts a precursor within the cell into a secondary messenger, which closes potassium ion channels. Also, this secondary messenger can stimulate the endoplasmic reticulum to release Ca2+ which contributes to depolarization. This leads to a build-up of potassium ions in the cell, depolarization, and neurotransmitter release. It is also possible for some bitter tastants to interact directly with the G protein, because of a structural similarity to the relevant GPCR.
Savoriness
Savoriness, or umami, is an appetitive taste. It can be tasted in soy sauce, meat, dashi and consomme. Umami, a loanword from Japanese meaning "good flavor" or "good taste", which is similar to the word "savory" that comes from the French for "tasty". is considered fundamental to many East Asian cuisines, such as Japanese cuisine. It dates back to the use of fermented fish sauce: garum in ancient Rome and or in ancient China.
Umami was first studied in 1907 by Ikeda isolating dashi taste, which he identified as the chemical monosodium glutamate (MSG). MSG is a sodium salt that produces a strong savory taste, especially combined with foods rich in nucleotides such as meats, fish, nuts, and mushrooms.
Some savory taste buds respond specifically to glutamate in the same way that "sweet" ones respond to sugar. Glutamate binds to a variant of G protein coupled glutamate receptors. L-glutamate may bond to a type of GPCR known as a metabotropic glutamate receptor (mGluR4) which causes the G-protein complex to activate the sensation of umami.
Perceptual independence from salty and sweet taste
There are doubts regarding whether umami is different from salty taste, as standalone glutamate(glutamic acid) without table salt ions(Na+), is perceived as sour, salt taste blockers reduce discrimination between monosodium glutamate and sucrose in rodents, since sweet and umami tastes share a taste receptor subunit; and part of the human population cannot tell apart umami from salty.
If umami doesn't have perceptual independence, it could be classified with other tastes like fat, carbohydrate, metallic, and calcium, which can be perceived at high concentrations but may not offer a prominent taste experience.
Measuring relative tastes
Measuring the degree to which a substance presents one basic taste can be achieved in a subjective way by comparing its taste to a reference substance.
Sweetness is subjectively measured by comparing the threshold values, or level at which the presence of a dilute substance can be detected by a human taster, of different sweet substances. Substances are usually measured relative to sucrose, which is usually given an arbitrary index of 1 or 100. Rebaudioside A is 100 times sweeter than sucrose; fructose is about 1.4 times sweeter; glucose, a sugar found in honey and vegetables, is about three-quarters as sweet; and lactose, a milk sugar, is one-half as sweet.
The sourness of a substance can be rated by comparing it to very dilute hydrochloric acid (HCl).
Relative saltiness can be rated by comparison to a dilute salt solution.
Quinine, a bitter medicinal found in tonic water, can be used to subjectively rate the bitterness of a substance. Units of dilute quinine hydrochloride (1 g in 2000 mL of water) can be used to measure the threshold bitterness concentration, the level at which the presence of a dilute bitter substance can be detected by a human taster, of other compounds. More formal chemical analysis, while possible, is difficult.
There may not be an absolute measure for pungency, though there are tests for measuring the subjective presence of a given pungent substance in food, such as the Scoville scale for capsaicine in peppers or the Pyruvate scale for pyruvates in garlics and onions.
Functional structure
Taste is a form of chemoreception which occurs in the specialised taste receptors in the mouth. To date, there are five different types of taste these receptors can detect which are recognized: salt, sweet, sour, bitter, and umami. Each type of receptor has a different manner of sensory transduction: that is, of detecting the presence of a certain compound and starting an action potential which alerts the brain. It is a matter of debate whether each taste cell is tuned to one specific tastant or to several; Smith and Margolskee claim that "gustatory neurons typically respond to more than one kind of stimulus, [a]lthough each neuron responds most strongly to one tastant". Researchers believe that the brain interprets complex tastes by examining patterns from a large set of neuron responses. This enables the body to make "keep or spit out" decisions when there is more than one tastant present. "No single neuron type alone is capable of discriminating among stimuli or different qualities, because a given cell can respond the same way to disparate stimuli." As well, serotonin is thought to act as an intermediary hormone which communicates with taste cells within a taste bud, mediating the signals being sent to the brain. Receptor molecules are found on the top of microvilli of the taste cells.
Sweetness
Sweetness is produced by the presence of sugars, some proteins, and other substances such as alcohols like anethol, glycerol and propylene glycol, saponins such as glycyrrhizin, artificial sweeteners (organic compounds with a variety of structures), and lead compounds such as lead acetate. It is often connected to aldehydes and ketones, which contain a carbonyl group. Many foods can be perceived as sweet regardless of their actual sugar content. For example, some plants such as liquorice, anise or stevia can be used as sweeteners. Rebaudioside A is a steviol glycoside coming from stevia that is 200 times sweeter than sugar. Lead acetate and other lead compounds were used as sweeteners, mostly for wine, until lead poisoning became known. Romans used to deliberately boil the must inside of lead vessels to make a sweeter wine.
Sweetness is detected by a variety of G protein-coupled receptors coupled to a G protein that acts as an intermediary in the communication between taste bud and brain, gustducin. These receptors are T1R2+3 (heterodimer) and T1R3 (homodimer), which account for sweet sensing in humans and other animals.
Saltiness
Saltiness is a taste produced best by the presence of cations (such as , or ) and is directly detected by cation influx into glial like cells via leak channels causing depolarisation of the cell.
Other monovalent cations, e.g., ammonium, , and divalent cations of the alkali earth metal group of the periodic table, e.g., calcium, , ions, in general, elicit a bitter rather than a salty taste even though they, too, can pass directly through ion channels in the tongue.
Sourness
Sourness is acidity, and, like salt, it is a taste sensed using ion channels. Undissociated acid diffuses across the plasma membrane of a presynaptic cell, where it dissociates in accordance with Le Chatelier's principle. The protons that are released then block potassium channels, which depolarise the cell and cause calcium influx. In addition, the taste receptor PKD2L1 has been found to be involved in tasting sour.
Bitterness
Research has shown that TAS2Rs (taste receptors, type 2, also known as T2Rs) such as TAS2R38 are responsible for the ability to taste bitter substances in vertebrates. They are identified not only by their ability to taste certain bitter ligands, but also by the morphology of the receptor itself (surface bound, monomeric).
Savoriness
The amino acid glutamic acid is responsible for savoriness, but some nucleotides (inosinic acid and guanylic acid) can act as complements, enhancing the taste.
Glutamic acid binds to a variant of the G protein-coupled receptor, producing a savory taste.
Further sensations and transmission
The tongue can also feel other sensations not generally included in the basic tastes. These are largely detected by the somatosensory system. In humans, the sense of taste is conveyed via three of the twelve cranial nerves. The facial nerve (VII) carries taste sensations from the anterior two thirds of the tongue, the glossopharyngeal nerve (IX) carries taste sensations from the posterior one third of the tongue while a branch of the vagus nerve (X) carries some taste sensations from the back of the oral cavity.
The trigeminal nerve (cranial nerve V) provides information concerning the general texture of food as well as the taste-related sensations of peppery or hot (from spices).
Pungency (also spiciness or hotness)
Substances such as ethanol and capsaicin cause a burning sensation by inducing a trigeminal nerve reaction together with normal taste reception. The sensation of heat is caused by the food's activating nerves that express TRPV1 and TRPA1 receptors. Some such plant-derived compounds that provide this sensation are capsaicin from chili peppers, piperine from black pepper, gingerol from ginger root and allyl isothiocyanate from horseradish. The piquant ("hot" or "spicy") sensation provided by such foods and spices plays an important role in a diverse range of cuisines across the world—especially in equatorial and sub-tropical climates, such as Ethiopian, Peruvian, Hungarian, Indian, Korean, Indonesian, Lao, Malaysian, Mexican, New Mexican, Pakistani, Singaporean, Southwest Chinese (including Sichuan cuisine), Vietnamese, and Thai cuisines.
This particular sensation, called chemesthesis, is not a taste in the technical sense, because the sensation does not arise from taste buds, and a different set of nerve fibers carry it to the brain. Foods like chili peppers activate nerve fibers directly; the sensation interpreted as "hot" results from the stimulation of somatosensory (pain/temperature) fibers on the tongue. Many parts of the body with exposed membranes but no taste sensors (such as the nasal cavity, under the fingernails, surface of the eye or a wound) produce a similar sensation of heat when exposed to hotness agents.
Coolness
Some substances activate cold trigeminal receptors even when not at low temperatures. This "fresh" or "minty" sensation can be tasted in peppermint and spearmint and is triggered by substances such as menthol, anethol, ethanol, and camphor. Caused by activation of the same mechanism that signals cold, TRPM8 ion channels on nerve cells, unlike the actual change in temperature described for sugar substitutes, this coolness is only a perceived phenomenon.
Numbness
Both Chinese and Batak Toba cooking include the idea of 麻 (má or mati rasa), a tingling numbness caused by spices such as Sichuan pepper. The cuisines of Sichuan province in China and of the Indonesian province of North Sumatra often combine this with chili pepper to produce a 麻辣 málà, "numbing-and-hot", or "mati rasa" flavor.
Typical in northern Brazilian cuisine, jambu is an herb used in dishes like tacacá.
These sensations, although not taste, fall into a category of chemesthesis.
Astringency
Some foods, such as unripe fruits, contain tannins or calcium oxalate that cause an astringent or puckering sensation of the mucous membrane of the mouth. Examples include tea, red wine, or rhubarb. Other terms for the astringent sensation are "dry", "rough", "harsh" (especially for wine), "tart" (normally referring to sourness), "rubbery", "hard" or "styptic".
Metallicness
A metallic taste may be caused by food and drink, certain medicines or amalgam dental fillings. It is generally considered an off flavor when present in food and drink. A metallic taste may be caused by galvanic reactions in the mouth. In the case where it is caused by dental work, the dissimilar metals used may produce a measurable current. Some artificial sweeteners are perceived to have a metallic taste, which is detected by the TRPV1 receptors. Many people consider blood to have a metallic taste. A metallic taste in the mouth is also a symptom of various medical conditions, in which case it may be classified under the symptoms dysgeusia or parageusia, referring to distortions of the sense of taste, and can be caused by medication, including saquinavir, zonisamide, and various kinds of chemotherapy, as well as occupational hazards, such as working with pesticides.
Fat taste
Recent research reveals a potential taste receptor called the CD36 receptor. CD36 was targeted as a possible lipid taste receptor because it binds to fat molecules (more specifically, long-chain fatty acids), and it has been localized to taste bud cells (specifically, the circumvallate and foliate papillae). There is a debate over whether we can truly taste fats, and supporters of human ability to taste free fatty acids (FFAs) have based the argument on a few main points: there is an evolutionary advantage to oral fat detection; a potential fat receptor has been located on taste bud cells; fatty acids evoke specific responses that activate gustatory neurons, similar to other currently accepted tastes; and, there is a physiological response to the presence of oral fat. Although CD36 has been studied primarily in mice, research examining human subjects' ability to taste fats found that those with high levels of CD36 expression were more sensitive to tasting fat than were those with low levels of CD36 expression; this study points to a clear association between CD36 receptor quantity and the ability to taste fat.
Other possible fat taste receptors have been identified. G protein-coupled receptors free fatty acid receptor 4 (also termed GPR120) and to a much lesser extent Free fatty acid receptor 1 (also termed GPR40) have been linked to fat taste, because their absence resulted in reduced preference to two types of fatty acid (linoleic acid and oleic acid), as well as decreased neuronal response to oral fatty acids.
Monovalent cation channel TRPM5 has been implicated in fat taste as well, but it is thought to be involved primarily in downstream processing of the taste rather than primary reception, as it is with other tastes such as bitter, sweet, and savory.
Proposed alternate names to fat taste include oleogustus and pinguis, although these terms are not widely accepted. The main form of fat that is commonly ingested is triglycerides, which are composed of three fatty acids bound together. In this state, triglycerides are able to give fatty foods unique textures that are often described as creaminess. But this texture is not an actual taste. It is only during ingestion that the fatty acids that make up triglycerides are hydrolysed into fatty acids via lipases. The taste is commonly related to other, more negative, tastes such as bitter and sour due to how unpleasant the taste is for humans. Richard Mattes, a co-author of the study, explained that low concentrations of these fatty acids can create an overall better flavor in a food, much like how small uses of bitterness can make certain foods more rounded. A high concentration of fatty acids in certain foods is generally considered inedible. To demonstrate that individuals can distinguish fat taste from other tastes, the researchers separated volunteers into groups and had them try samples that also contained the other basic tastes. Volunteers were able to separate the taste of fatty acids into their own category, with some overlap with savory samples, which the researchers hypothesized was due to poor familiarity with both. The researchers note that the usual "creaminess and viscosity we associate with fatty foods is largely due to triglycerides", unrelated to the taste; while the actual taste of fatty acids is not pleasant. Mattes described the taste as "more of a warning system" that a certain food should not be eaten.
There are few regularly consumed foods rich in fat taste, due to the negative flavor that is evoked in large quantities. Foods whose flavor to which fat taste makes a small contribution include olive oil and fresh butter, along with various kinds of vegetable and nut oils.
Heartiness
Kokumi (, Japanese: from ) is translated as "heartiness", "full flavor" or "rich" and describes compounds in food that do not have their own taste, but enhance the characteristics when combined.
Alongside the five basic tastes of sweet, sour, salt, bitter and savory, kokumi has been described as something that may enhance the other five tastes by magnifying and lengthening the other tastes, or "mouthfulness". Garlic is a common ingredient to add flavor used to help define the characteristic kokumi flavors.
Calcium-sensing receptors (CaSR) are receptors for kokumi substances which, applied around taste pores, induce an increase in the intracellular Ca concentration in a subset of cells. This subset of CaSR-expressing taste cells are independent from the influenced basic taste receptor cells. CaSR agonists directly activate the CaSR on the surface of taste cells and integrated in the brain via the central nervous system. A basal level of calcium, corresponding to the physiological concentration, is necessary for activation of the CaSR to develop the kokumi sensation.
Calcium
The distinctive taste of chalk has been identified as the calcium component of that substance. In 2008, geneticists discovered a calcium receptor on the tongues of mice. The CaSR receptor is commonly found in the gastrointestinal tract, kidneys, and brain. Along with the "sweet" T1R3 receptor, the CaSR receptor can detect calcium as a taste. Whether the perception exists or not in humans is unknown.
Temperature
Temperature can be an essential element of the taste experience. Heat can accentuate some flavors and decrease others by varying the density and phase equilibrium of a substance. Food and drink that—in a given culture—is traditionally served hot is often considered distasteful if cold, and vice versa. For example, alcoholic beverages, with a few exceptions, are usually thought best when served at room temperature or chilled to varying degrees, but soups—again, with exceptions—are usually only eaten hot. A cultural example are soft drinks. In North America it is almost always preferred cold, regardless of season.
Starchiness
A 2016 study suggested that humans can taste starch (specifically, a glucose oligomer) independently of other tastes such as sweetness, without suggesting an associated chemical receptor.
Nerve supply and neural connections
The glossopharyngeal nerve innervates a third of the tongue including the circumvallate papillae. The facial nerve innervates the other two thirds of the tongue and the cheek via the chorda tympani.
The pterygopalatine ganglia are ganglia (one on each side) of the soft palate. The greater petrosal, lesser palatine and zygomatic nerves all synapse here. The greater petrosal carries soft palate taste signals to the facial nerve. The lesser palatine sends signals to the nasal cavity, which is why spicy foods cause nasal drip. The zygomatic sends signals to the lacrimal nerve that activate the lacrimal gland, which is the reason that spicy foods can cause tears. Both the lesser palatine and the zygomatic are maxillary nerves (from the trigeminal nerve).
The special visceral afferents of the vagus nerve carry taste from the epiglottal region of the tongue.
The lingual nerve (trigeminal, not shown in diagram) is deeply interconnected with the chorda tympani in that it provides all other sensory info from the anterior ⅔ of the tongue. This info is processed separately (nearby) in the rostral lateral subdivision of the nucleus of the solitary tract (NST).
The NST receives input from the amygdala (regulates oculomotor nuclei output), bed nuclei of stria terminalis, hypothalamus, and prefrontal cortex. The NST is the topographical map that processes gustatory and sensory (temp, texture, etc.) info.
The reticular formation (includes Raphe nuclei responsible for serotonin production) is signaled to release serotonin during and after a meal to suppress appetite. Similarly, salivary nuclei are signaled to decrease saliva secretion.
Hypoglossal and thalamic connections aid in oral-related movements.
Hypothalamus connections hormonally regulate hunger and the digestive system.
Substantia innominata connects the thalamus, temporal lobe, and insula.
Edinger-Westphal nucleus reacts to taste stimuli by dilating and constricting the pupils.
Spinal ganglia are involved in movement.
The frontal operculum is speculated to be the memory and association hub for taste.
The insula cortex aids in swallowing and gastric motility.
Taste in insects
Insects taste using small hair-like structures called taste sensilla, specialized sensory organs located on various body parts such as the mouthparts, legs, and wings. These sensilla contain gustatory receptor neurons (GRNs) sensitive to a wide range of chemical stimuli.
Insects respond to sugar, bitter, acid, and salt tastes. However, their taste spectrum extends to include water, fatty acids, metals, carbonation, RNA, ATP, and pheromones. Detecting these substances is vital for behaviors like feeding, mating, and oviposition.
Invertebrates' ability to taste these compounds is fundamental to their survival and provides insights into the evolution of sensory systems. This knowledge is crucial for understanding insect behavior and has applications in pest control and pollination biology.
Other concepts
Supertasters
A supertaster is a person whose sense of taste is significantly more sensitive than most. The cause of this heightened response is likely, at least in part, due to an increased number of fungiform papillae. Studies have shown that supertasters require less fat and sugar in their food to get the same satisfying effects. These people tend to consume more salt than others. This is due to their heightened sense of the taste of bitterness, and the presence of salt drowns out the taste of bitterness.
Aftertaste
Aftertastes arise after food has been swallowed. An aftertaste can differ from the food it follows. Medicines and tablets may also have a lingering aftertaste, as they can contain certain artificial flavor compounds, such as aspartame (artificial sweetener).
Acquired taste
An acquired taste often refers to an appreciation for a food or beverage that is unlikely to be enjoyed by a person who has not had substantial exposure to it, usually because of some unfamiliar aspect of the food or beverage, including bitterness, a strong or strange odor, taste, or appearance.
Clinical significance
Patients with Addison's disease, pituitary insufficiency, or cystic fibrosis sometimes have a hyper-sensitivity to the five primary tastes.
Disorders of taste
ageusia (complete loss of taste)
hypogeusia (reduced sense of taste)
dysgeusia (distortion in sense of taste)
hypergeusia (abnormally heightened sense of taste)
Viruses can also cause loss of taste. About 50% of patients with SARS-CoV-2 (causing COVID-19) experience some type of disorder associated with their sense of smell or taste, including ageusia and dysgeusia. SARS-CoV-1, MERS-CoV and even the flu (influenza virus) can also disrupt olfaction.
History
In the West, Aristotle postulated in BC that the two most basic tastes were sweet and bitter. He was one of the first persons to develop a list of basic tastes.
Research
The receptors for the basic tastes of bitter, sweet and savory have been identified. They are G protein-coupled receptors. The cells that detect sourness have been identified as a subpopulation that express the protein PKD2L1, and The responses are mediated by an influx of protons into the cells. As of 2019, molecular mechanisms for each taste appear to be different, although all taste perception relies on activation of P2X purinoreceptors on sensory nerves.
| Biology and health sciences | Nervous system | null |
21290714 | https://en.wikipedia.org/wiki/Proprioception | Proprioception | Proprioception ( ) is the sense of self-movement, force, and body position.
Proprioception is mediated by proprioceptors, a type of sensory receptor, located within muscles, tendons, and joints. Most animals possess multiple subtypes of proprioceptors, which detect distinct kinesthetic parameters, such as joint position, movement, and load. Although all mobile animals possess proprioceptors, the structure of the sensory organs can vary across species.
Proprioceptive signals are transmitted to the central nervous system, where they are integrated with information from other sensory systems, such as the visual system and the vestibular system, to create an overall representation of body position, movement, and acceleration. In many animals, sensory feedback from proprioceptors is essential for stabilizing body posture and coordinating body movement.
System overview
In vertebrates, limb movement and velocity (muscle length and the rate of change) are encoded by one group of sensory neurons (type Ia sensory fiber) and another type encode static muscle length (group II neurons). These two types of sensory neurons compose muscle spindles. There is a similar division of encoding in invertebrates; different subgroups of neurons of the chordotonal organ encode limb position and velocity.
To determine the load on a limb, vertebrates use sensory neurons in the Golgi tendon organs: type Ib afferents. These proprioceptors are activated at given muscle forces, which indicate the resistance that muscle is experiencing. Similarly, invertebrates have a mechanism to determine limb load: the campaniform sensilla. These proprioceptors are active when a limb experiences resistance.
A third role for proprioceptors is to determine when a joint is at a specific position. In vertebrates, this is accomplished by Ruffini endings and Pacinian corpuscles. These proprioceptors are activated when the joint is at a threshold position, usually at the extremes of joint position. Invertebrates use hair plates to accomplish this; a field of bristles located within joints that detects the relative movement of limb segments through the deflection of the associated cuticular hairs.
Reflexes
The sense of proprioception is ubiquitous across mobile animals and is essential for the motor coordination of the body. Proprioceptors can form reflex circuits with motor neurons to provide rapid feedback about body and limb position. These mechanosensation circuits are important for flexibly maintaining posture and balance, especially during locomotion. For example, consider the stretch reflex, in which stretch across a muscle is detected by a sensory receptor (e.g., muscle spindle, chordotonal neurons), which activates a motor neuron to induce muscle contraction and oppose the stretch. During locomotion, sensory neurons can reverse their activity when stretched, to promote rather than oppose movement.
Conscious and nonconscious
In humans, a distinction is made between conscious proprioception and nonconscious proprioception:
Conscious proprioception is communicated by the dorsal column-medial lemniscus pathway to the cerebrum.
Nonconscious proprioception is communicated primarily via the dorsal spinocerebellar tract and ventral spinocerebellar tract, to the cerebellum.
A nonconscious reaction is seen in the human proprioceptive reflex, or righting reflex—in the event that the body tilts in any direction, the person will cock their head back to level the eyes against the horizon. This is seen even in infants as soon as they gain control of their neck muscles. This control comes from the cerebellum, the part of the brain affecting balance.
Mechanisms
Proprioception is mediated by mechanically sensitive proprioceptor neurons distributed throughout an animal's body. Most vertebrates possess three basic types of proprioceptors: muscle spindles, which are embedded in skeletal muscles, Golgi tendon organs, which lie at the interface of muscles and tendons, and joint receptors, which are low-threshold mechanoreceptors embedded in joint capsules. Many invertebrates, such as insects, also possess three basic proprioceptor types with analogous functional properties: chordotonal neurons, campaniform sensilla, and hair plates.
The initiation of proprioception is the activation of a proprioceptor in the periphery. The proprioceptive sense is believed to be composed of information from sensory neurons located in the inner ear (motion and orientation) and in the stretch receptors located in the muscles and the joint-supporting ligaments (stance). There are specific nerve receptors for this form of perception termed "proprioceptors", just as there are specific receptors for pressure, light, temperature, sound, and other sensory experiences. Proprioceptors are sometimes known as adequate stimuli receptors.
Members of the transient receptor potential family of ion channels have been found to be important for proprioception in fruit flies, nematode worms, African clawed frogs, and zebrafish. PIEZO2, a nonselective cation channel, has been shown to underlie the mechanosensitivity of proprioceptors in mice. Humans with loss-of-function mutations in the PIEZO2 gene exhibit specific deficits in joint proprioception, as well as vibration and touch discrimination, suggesting that the PIEZO2 channel is essential for mechanosensitivity in some proprioceptors and low-threshold mechanoreceptors.
Although it was known that finger kinesthesia relies on skin sensation, recent research has found that kinesthesia-based haptic perception relies strongly on the forces experienced during touch. This research allows the creation of "virtual", illusory haptic shapes with different perceived qualities.
Anatomy
Proprioception of the head stems from the muscles innervated by the trigeminal nerve, where the general somatic afferent fibers pass without synapsing in the trigeminal ganglion (first-order sensory neuron), reaching the mesencephalic tract and the mesencephalic nucleus of trigeminal nerve. Proprioception of limbs often occurs due to receptors in connective tissue near joints.
Function
Stability
An important role for proprioception is to allow an animal to stabilize itself against perturbations. For instance, for a person to walk or stand upright, they must continuously monitor their posture and adjust muscle activity as needed to provide balance. Similarly, when walking on unfamiliar terrain, or even tripping, the person must adjust the output of their muscles quickly based on estimated limb position and velocity. Proprioceptor reflex circuits are thought to play an important role to allow fast and unconscious execution of these behaviors, To make control of these behaviors efficient, proprioceptors are also thought to regulate reciprocal inhibition in muscles, leading to agonist-antagonist muscle pairs.
Planning and refining movements
When planning complex movements such as reaching or grooming, an animal must consider the current position and velocity of its limb and use that information to adjust dynamics to target a final position. If the animal's estimate of its limb's initial position is wrong, then a deficiency in the movement can result. Furthermore, proprioception is crucial in refining the movement if it deviates from the trajectory.
Development
In adult fruit flies, each proprioceptor class arises from a specific cell lineage (i.e. each chordotonal neuron is from the chordotonal neuron lineage, although multiple lineages give rise to sensory bristles). After the last cell division, proprioceptors send out axons toward the central nervous system and are guided by hormonal gradients to reach stereotyped synapses.
The mechanisms underlying axon guidance are similar across invertebrates and vertebrates.
In mammals with longer gestation periods, muscle spindles are fully formed at birth. Muscle spindles continue to grow throughout post-natal development as muscles grow.
Mathematical models
Proprioceptors transfer the mechanical state of the body into patterns of neural activity. This transfer can be modeled mathematically, for example to better understand the internal workings of a proprioceptor or to provide more realistic feedback in neuromechanical simulations.
Various proprioceptor models of complexity have been developed. They range from simple phenomenological models to complex structural models, in which the mathematical elements correspond to anatomical features of the proprioceptor. The focus has been on muscle spindles, but Golgi tendon organs and insects' hair plates have been modeled too.
Muscle spindles
Poppele and Bowman used linear system theory to model mammalian muscle spindles Ia and II afferents. They obtained a set of de-afferented muscle spindles, measured their response to a series of sinusoidal and step function stretches, and fit a transfer function to the spike rate. They found that the following Laplace transfer function describes the firing rate responses of the primary sensory fibers for a change in length:
The following equation describes the response of secondary sensory fibers:
More recently, Blum et al. showed that the muscle spindle firing rate is modeled better as tracking the force of the muscle, rather than the length. Furthermore, muscle spindle firing rates show history dependence which cannot be modeled by a linear time-invariant system model.
Golgi tendon organs
Houk and Simon provided one of the first mathematical models of a Golgi tendon organ receptor, modeling the firing rate of the receptor as a function of the muscle tension force. Just as for muscle spindles, they find that, as the receptors respond linearly to sine waves of different frequencies and has little variance in response over time to the same stimulus, Golgi tendon organ receptors may be modeled as linear time-invariant systems. Specifically, they find that the firing rate of a Golgi tendon organ receptor may be modeled as a sum of 3 decaying exponentials:
where is the firing rate and is a step function of force.
The corresponding Laplace transfer function for this system is:
For a soleus receptor, Houk and Simon obtain average values of K=57 pulses/sec/kg, A=0.31, a=0.22 sec−1, B=0.4, b=2.17 sec−1, C=2.5, c=36 sec−1 .
When modeling a stretch reflex, Lin and Crago improved upon this model by adding a logarithmic nonlinearity before the Houk and Simon model and a threshold nonlinearity after.
Impairment
Chronic
Proprioception, a sense vital for rapid and proper body coordination, can be permanently lost or impaired as a result of genetic conditions, disease, viral infections, and injuries. For instance, patients with joint hypermobility or Ehlers–Danlos syndromes, genetic conditions that result in weak connective tissue throughout the body, have chronic impairments to proprioception. Autism spectrum disorder and Parkinson's disease can also cause chronic disorder of proprioception. In regards to Parkinson's disease, it remains unclear whether the proprioceptive-related decline in motor function occurs due to disrupted proprioceptors in the periphery or signaling in the spinal cord or brain.
In rare cases, viral infections result in a loss of proprioception. Ian Waterman and Charles Freed are two such people that lost their sense of proprioception from the neck down from supposed viral infections (i.e. gastric flu and a rare viral infection). After losing their sense of proprioception, Ian and Charles could move their lower body, but could not coordinate their movements. However, both individuals regained some control of their limbs and body by consciously planning their movements and relying solely on visual feedback. Interestingly, both individuals can still sense pain and temperature, indicating that they specifically lost proprioceptive feedback, but not tactile and nociceptive feedback. The impact of losing the sense of proprioception on daily life is perfectly illustrated when Ian Waterman stated, "What is an active brain without mobility".
Proprioception is also permanently lost in people who lose a limb or body part through injury or amputation. After the removal of a limb, people may have a confused sense of that limb's existence on their body, known as phantom limb syndrome. Phantom sensations can occur as passive proprioceptive sensations of the limb's presence, or more active sensations such as perceived movement, pressure, pain, itching, or temperature. There are a variety of theories concerning the etiology of phantom limb sensations and experience. One is the concept of "proprioceptive memory", which argues that the brain retains a memory of specific limb positions and that after amputation there is a conflict between the visual system, which actually sees that the limb is missing, and the memory system which remembers the limb as a functioning part of the body. Phantom sensations and phantom pain may also occur after the removal of body parts other than the limbs, such as after amputation of the breast, extraction of a tooth (phantom tooth pain), or removal of an eye (phantom eye syndrome).
There is a decline in the sense of proprioception with ageing. This can often result in chronic lower back pain, and be the cause of falls in the elderly.
Acute
Proprioception is occasionally impaired spontaneously, especially when one is tired. Similar effects can be felt during the hypnagogic state of consciousness, during the onset of sleep. One's body may feel too large or too small, or parts of the body may feel distorted in size. Similar effects can sometimes occur during epilepsy or migraine auras. These effects are presumed to arise from abnormal stimulation of the part of the parietal cortex of the brain involved with integrating information from different parts of the body. Proprioceptive illusions can also be induced, such as the "Pinocchio illusion", the illusion that one's nose is growing longer.
Temporary impairment of proprioception has also been known to occur from an overdose of vitamin B6 (pyridoxine and pyridoxamine). This is due to a reversible neuropathy. Most of the impaired function returns to normal shortly after the amount of the vitamin in the body returns to a level that is closer to that of the physiological norm. Impairment can also be caused by cytotoxic factors such as chemotherapy.
It has been proposed that even common tinnitus and the attendant hearing frequency-gaps masked by the perceived sounds may cause erroneous proprioceptive information to the balance and comprehension centers of the brain, precipitating mild confusion.
Temporary loss or impairment of proprioception may happen periodically during growth, mostly during adolescence. Growth that might also influence this would be large increases or drops in bodyweight/size due to fluctuations of fat (liposuction, rapid fat loss or gain) and/or muscle content (bodybuilding, anabolic steroids, catabolisis/starvation). It can also occur in those that gain new levels of flexibility, stretching, and contortion. A limb's being in a new range of motion never experienced (or at least, not for a long time since youth perhaps) can disrupt one's sense of location of that limb. Possible experiences include suddenly feeling that feet or legs are missing from one's mental self-image; needing to look down at one's limbs to be sure they are still there; and falling down while walking, especially when attention is focused upon something other than the act of walking.
Diagnosis
Impaired proprioception may be diagnosed through a series of tests, each focusing on a different functional aspect of proprioception.
The Romberg's test is often used to assess balance. The subject must stand with feet together and eyes closed without support for 30 seconds. If the subject loses balance and falls, it is an indicator for impaired proprioception.
For evaluating proprioception's contribution to motor control, a common protocol is joint position matching. The patient is blindfolded while a joint is moved to a specific angle for a given period of time and then returned to neutral. The subject is then asked to move the joint back to the specified angle. Recent investigations have shown that hand dominance, participant age, active versus passive matching, and presentation time of the angle can all affect performance on joint position matching tasks.
For passive sensing of joint angles, recent studies have found that experiments to probe psychophysical thresholds produce more precise estimates of proprioceptive discrimination than the joint position matching task. In these experiments, the subject holds on to an object (such as an armrest) that moves and stops at different positions. The subject must discriminate whether one position is closer to the body than another. From the subject's choices, the tester may determine the subject's discrimination thresholds.
Proprioception is tested by American police officers using the field sobriety testing to check for alcohol intoxication. The subject is required to touch his or her nose with eyes closed; people with normal proprioception may make an error of no more than , while people with impaired proprioception (a symptom of moderate to severe alcohol intoxication) fail this test due to difficulty locating their limbs in space relative to their noses.
Training
Proprioception is what allows someone to learn to walk in complete darkness without losing balance. During the learning of any new skill, sport, or art, it is usually necessary to become familiar with some proprioceptive tasks specific to that activity. Without the appropriate integration of proprioceptive input, an artist would not be able to brush paint onto a canvas without looking at the hand as it moved the brush over the canvas; it would be impossible to drive an automobile because a motorist would not be able to steer or use the pedals while looking at the road ahead; a person could not touch type or perform ballet; and people would not even be able to walk without watching where they put their feet.
Oliver Sacks reported the case of a young woman who lost her proprioception due to a viral infection of her spinal cord. At first she could not move properly at all or even control her tone of voice (as voice modulation is primarily proprioceptive). Later she relearned by using her sight (watching her feet) and inner ear only for movement while using hearing to judge voice modulation. She eventually acquired a stiff and slow movement and nearly normal speech, which is believed to be the best possible in the absence of this sense. She could not judge effort involved in picking up objects and would grip them painfully to be sure she did not drop them.
The proprioceptive sense can be sharpened through study of many disciplines. Juggling trains reaction time, spatial location, and efficient movement. Standing on a wobble board or balance board is often used to retrain or increase proprioceptive abilities, particularly as physical therapy for ankle or knee injuries. Slacklining is another method to increase proprioception.
Standing on one leg (stork standing) and various other body-position challenges are also used in such disciplines as yoga, Wing Chun and tai chi. The vestibular system of the inner ear, vision and proprioception are the main three requirements for balance. Moreover, there are specific devices designed for proprioception training, such as the exercise ball, which works on balancing the abdominal and back muscles.
History of study
In 1557, the position-movement sensation was described by Julius Caesar Scaliger as a "sense of locomotion".
In 1826, Charles Bell expounded the idea of a "muscle sense", which is credited as one of the first descriptions of physiologic feedback mechanisms. Bell's idea was that commands are carried from the brain to the muscles, and that reports on the muscle's condition would be sent in the reverse direction.
In 1847, the London neurologist Robert Todd highlighted important differences in the anterolateral and posterior columns of the spinal cord, and suggested that the latter were involved in the coordination of movement and balance.
At around the same time, Moritz Heinrich Romberg, a Berlin neurologist, was describing unsteadiness made worse by eye closure or darkness, now known as the eponymous Romberg's sign, once synonymous with tabes dorsalis, that became recognised as common to all proprioceptive disorders of the legs.
In 1880, Henry Charlton Bastian suggested "kinaesthesia" instead of "muscle sense" on the basis that some of the afferent information (back to the brain) comes from other structures, including tendons, joints, and skin.
In 1889, Alfred Goldscheider suggested a classification of kinaesthesia into three types: muscle, tendon, and articular sensitivity.
In 1906, the term proprio-ception (and also intero-ception and extero-ception) is attested in a publication by Charles Scott Sherrington involving receptors. He explains the terminology as follows:
Today, the "exteroceptors" are the organs that provide information originating outside the body, such as the eyes, ears, mouth, and skin. The interoceptors provide information about the internal organs, and the "proprioceptors" provide information about movement derived from muscular, tendon, and articular sources. Using Sherrington's system, physiologists and anatomists search for specialised nerve endings that transmit mechanical data on joint capsule, tendon and muscle tension (such as Golgi tendon organs and muscle spindles), which play a large role in proprioception.
Primary endings of muscle spindles "respond to the size of a muscle length change and its speed" and "contribute both to the sense of limb position and movement". Secondary endings of muscle spindles detect changes in muscle length, and thus supply information regarding only the sense of position. Essentially, muscle spindles are stretch receptors. It has been accepted that cutaneous receptors also contribute directly to proprioception by providing "accurate perceptual information about joint position and movement", and this knowledge is combined with information from the muscle spindles.
Etymology
Proprioception is from Latin proprius, meaning "one's own", "individual", and capio, capere, to take or grasp. Thus to grasp one's own position in space, including the position of the limbs in relation to each other and the body as a whole.
The word kinesthesia or kinæsthesia (kinesthetic sense) refers to movement sense, but has been used inconsistently to refer either to proprioception alone or to the brain's integration of proprioceptive and vestibular inputs. Kinesthesia is a modern medical term composed of elements from Greek; kinein "to set in motion; to move" (from PIE root *keie- "to set in motion") + aisthesis "perception, feeling" (from PIE root *au- "to perceive").
Plants and bacteria
Although they lack neurons, systems responding to stimuli (analogous to the sensory system in animals with a nervous system, which includes the proprioception) have also been described in some plants (angiosperms). Terrestrial plants control the orientation of their primary growth through the sensing of several vectorial stimuli such as the light gradient or the gravitational acceleration. This control has been called tropism. A quantitative study of shoot gravitropism demonstrated that, when a plant is tilted, it cannot recover a steady erected posture under the sole driving of the sensing of its angular deflection versus gravity. An additional control through the continuous sensing of its curvature by the organ and the subsequent driving an active straightening process are required. Being a sensing by the plant of the relative configuration of its parts, it has been called proprioception. This dual sensing and control by gravisensing and proprioception has been formalized into a unifying mathematical model simulating the complete driving of the gravitropic movement. This model has been validated on 11 species sampling the phylogeny of land angiosperms, and on organs of very contrasted sizes, ranging from the small germination of wheat (coleoptile) to the trunk of poplar trees.
Further studies have shown that the cellular mechanism of proprioception in plants involves myosin and actin, and seems to occur in specialized cells. Proprioception was then found to be involved in other tropisms and to be central also to the control of nutation.
The discovery of proprioception in plants has generated an interest in the popular science and generalist media. This is because this discovery questions a long-lasting a priori that we have on plants. In some cases this has led to a shift between proprioception and self-awareness or self-consciousness. There is no scientific ground for such a semantic shift. Indeed, even in animals, proprioception can be unconscious; so it is thought to be in plants.
Recent studies suggest that bacteria have control systems that may resemble proprioception.
| Biology and health sciences | Sensory nervous system | Biology |
21291483 | https://en.wikipedia.org/wiki/Windows%20NT | Windows NT | Windows NT is a proprietary graphical operating system produced by Microsoft as part of its Windows product line, the first version of which, Windows NT 3.1, was released on July 27, 1993. Originally made for the workstation, office, and server markets, the Windows NT line was made available to consumers with the release of Windows XP in 2001. The underlying technology of Windows NT continues to exist to this day with incremental changes and improvements, with the latest version of Windows based on Windows NT being Windows Server 2025 announced in 2024.
The name "Windows NT" originally denoted the major technological advancements that it had introduced to the Windows product line, including eliminating the 16-bit memory access limitations of earlier Windows releases such as Windows 3.1 and the Windows 9x series. Each Windows release built on this technology is considered to be based on, if not a revision of Windows NT, even though the Windows NT name itself has not been used in many other Windows releases since Windows NT 4.0 in 1996.
Windows NT provides many more features than other Windows releases, among them being support for multiprocessing, multi-user systems, a "pure" 32-bit kernel with 32-bit memory addressing, support for instruction sets other than x86, and many other system services such as Active Directory and more. Newer versions of Windows NT support 64-bit computing, with a 64-bit kernel and 64-bit memory addressing.
Product line
Windows NT is a group or family of products—like Windows is a group or family. Windows NT is a sub-grouping of Windows.
The first version of Windows NT, 3.1, was produced for workstation and server computers. It was commercially focused—and intended to complement consumer versions of Windows that were based on MS-DOS (including Windows 1.0 through Windows 3.1x). In 1996, Windows NT 4.0 was released, including the new shell from Windows 95.
Eventually, Microsoft incorporated the Windows NT technology into the Windows product line for personal computing and deprecated the Windows 9x family. Starting with Windows 2000, "NT" was removed from the product name yet is still in several low-level places in the system—including for a while as part of the product version.
Installing
Versions of Windows NT are installed using Windows Setup, which, starting with Windows Vista, uses the Windows Preinstallation Environment, which is a lightweight version of Windows NT made for deployment of the operating system.
Since Windows Vista, the Windows installation files, as well as the preinstallation environment used to install Windows, are stored in the Windows Imaging Format. It is possible to use the Deployment Image Servicing and Management (DISM) tool to install Windows from the command line and skip the GUI installer.
Naming
It has been suggested that Dave Cutler intended the initialism "WNT" as a play on VMS, incrementing each letter by one. However, the project was originally intended as a follow-on to OS/2 and was referred to as "NT OS/2" before receiving the Windows brand. One of the original NT developers, Mark Lucovsky, states that the name was taken from the original target processor—the Intel i860, code-named N10 ("N-Ten"). A 1991 video featuring Bill Gates and Microsoft products specifically says that "Windows NT stands for 'New Technology'". Seven years later in 1998, during a question-and-answer (Q&A) session, he then revealed that the letters were previously expanded to such but no longer carry any specific meaning. The letters were dropped from the names of releases from Windows 2000 and later, though Microsoft described that product as being "Built on NT Technology".
"NT" was a trademark of Northern Telecom (later Nortel), which Microsoft was forced to acknowledge on the product packaging.
Major features
One of the main purposes of NT is hardware and software portability. Various versions of NT family operating systems have been released for a variety of processor architectures, initially IA-32, MIPS, and DEC Alpha, with PowerPC, Itanium, x86-64 and ARM supported in later releases. An initial idea was to have a common code base with a custom Hardware Abstraction Layer (HAL) for each platform. However, support for MIPS, Alpha, and PowerPC was later dropped in Windows 2000. Broad software compatibility was initially achieved with support for several API "personalities", including Windows API, POSIX, and OS/2 APIs—the latter two were phased out starting with Windows XP. Partial MS-DOS and Windows 16-bit compatibility is achieved on IA-32 via an integrated DOS Virtual Machine—although this feature is not available on other architectures.
NT has supported per-object (file, function, and role) access control lists allowing a rich set of security permissions to be applied to systems and services. NT
has also supported Windows network protocols, inheriting the previous OS/2 LAN Manager networking, as well as TCP/IP networking (for which Microsoft used to implement a TCP/IP stack derived at first from a STREAMS-based stack from Spider Systems, then later rewritten in-house).
Windows NT 3.1 was the first version of Windows to use 32-bit flat virtual memory addressing on 32-bit processors. Its companion product, Windows 3.1, used segmented addressing and switches from 16-bit to 32-bit addressing in pages.
Windows NT 3.1 featured a core kernel providing a system API, running in supervisor mode (ring 0 in x86; referred to in Windows NT as "kernel mode" on all platforms), and a set of user-space environments with their own APIs which included the new Win32 environment, an OS/2 1.3 text-mode environment and a POSIX environment. The full preemptive multitasking kernel could interrupt running tasks to schedule other tasks, without relying on user programs to voluntarily give up control of the CPU, as in Windows 3.1 Windows applications (although MS-DOS applications were preemptively multitasked in Windows starting with Windows/386).
Notably, in Windows NT 3.x, several I/O driver subsystems, such as video and printing, were user-mode subsystems. In Windows NT 4.0, the video, server, and printer spooler subsystems were moved into kernel mode. Windows NT's first GUI was strongly influenced by (and programmatically compatible with) that from Windows 3.1; Windows NT 4.0's interface was redesigned to match that of the brand-new Windows 95, moving from the Program Manager to the Windows shell design.
NTFS, a journaled, secure file system, is a major feature of NT. Windows NT also allows for other installable file systems; NT can also be installed on FAT file systems, and versions 3.1, 3.5, and 3.51 could be installed HPFS file systems.
Windows NT introduced its own driver model, the Windows NT driver model, and is incompatible with older driver frameworks. With Windows 2000, the Windows NT driver model was enhanced to become the Windows Driver Model, which was first introduced with Windows 98, but was based on the NT driver model. Windows Vista added native support for the Windows Driver Foundation, which is also available for Windows XP, Windows Server 2003 and to an extent, Windows 2000.
Development
Microsoft decided to create a portable operating system, compatible with OS/2 and POSIX and supporting multiprocessing, in October 1988. When development started in November 1989, Windows NT was to be known as OS/2 3.0, the third version of the operating system developed jointly by Microsoft and IBM. To ensure portability, initial development was targeted at the Intel i860XR RISC processor, switching to the MIPS R3000 in late 1989, and then the Intel i386 in 1990. Microsoft also continued parallel development of the DOS-based and less resource-demanding Windows environment, resulting in the release of Windows 3.0 in May 1990.
Windows 3.0 was eventually so successful that Microsoft decided to change the primary application programming interface for the still unreleased NT OS/2 (as it was then known) from an extended OS/2 API to an extended Windows API. This decision caused tension between Microsoft and IBM and the collaboration ultimately fell apart.
IBM continued OS/2 development alone while Microsoft continued work on the newly renamed Windows NT. Though neither operating system would immediately be as popular as Microsoft's MS-DOS or Windows products, Windows NT would eventually be far more successful than OS/2.
Microsoft hired a group of developers from Digital Equipment Corporation led by Dave Cutler to build Windows NT, and many elements of the design reflect earlier DEC experience with Cutler's VMS, VAXELN and RSX-11, but also an unreleased object-based operating system developed by Cutler at Digital codenamed MICA. The team was joined by selected members of the disbanded OS/2 team, including Moshe Dunie.
Although NT was not an exact clone of Cutler's previous operating systems, DEC engineers almost immediately noticed the internal similarities. Parts of VAX/VMS Internals and Data Structures, published by Digital Press, accurately describe Windows NT internals using VMS terms. Furthermore, parts of the NT codebase's directory structure and filenames matched that of the MICA codebase. Instead of a lawsuit, Microsoft agreed to pay DEC $65–100 million, help market VMS, train Digital personnel on Windows NT, and continue Windows NT support for the DEC Alpha.
Windows NT and VMS memory management, processes, and scheduling are very similar. Windows NT's process management differs by implementing threading, which DEC did not implement until VMS 7.0 in 1995.
Like VMS, Windows NT's kernel mode code distinguishes between the "kernel", whose primary purpose is to implement processor- and architecture-dependent functions, and the "executive". This was designed as a modified microkernel, as the Windows NT kernel was influenced by the Mach microkernel developed by Richard Rashid at Carnegie Mellon University, but does not meet all of the criteria of a pure microkernel. Both the kernel and the executive are linked together into the single loaded module ntoskrnl.exe; from outside this module, there is little distinction between the kernel and the executive. Routines from each are directly accessible, as for example from kernel-mode device drivers.
API sets in the Windows NT family are implemented as subsystems atop the publicly undocumented "native" API; this allowed the late adoption of the Windows API (into the Win32 subsystem). Windows NT was one of the earliest operating systems to use UCS-2 and UTF-16 internally.
Architecture
Windows NT uses a layered design architecture that consists of two main components, user mode and kernel mode. Programs and subsystems in user mode are limited in terms of what system resources they have access to, while the kernel mode has unrestricted access to the system memory and external devices. Kernel mode in Windows NT has full access to the hardware and system resources of the computer. The Windows NT kernel is a hybrid kernel; the architecture comprises a simple kernel, hardware abstraction layer (HAL), drivers, and a range of services (collectively named Executive), which all exist in kernel mode.
The booting process of Windows NT begins with NTLDR in versions before Vista and the Windows Boot Manager in Vista and later. The boot loader is responsible for accessing the file system on the boot drive, starting the kernel, and loading boot-time device drivers into memory. Once all the boot and system drivers have been loaded, the kernel starts the Session Manager Subsystem. This process launches winlogon, which allows the user to login. Once the user is logged in File Explorer is started, loading the graphical user interface of Windows NT.
Programming language
Windows NT is written in C and C++, with a very small amount written in assembly language. C is mostly used for the kernel code while C++ is mostly used for user-mode code. Assembly language is avoided where possible because it would impede portability.
Releases
The following are the releases of Windows based on the Windows NT technology.
Windows NT 3.1 to 3.51 incorporated the Program Manager and File Manager from the Windows 3.1 series. Windows NT 4.0 onwards replaced those programs with Windows Explorer (including a taskbar and Start menu), which originally appeared in Windows 95.
The first release was given version number 3.1 to match the contemporary 16-bit Windows; magazines of that era claimed the number was also used to make that version seem more reliable than a ".0" release. Also the Novell IPX protocol was apparently licensed only to 3.1 versions of Windows software.
The NT version number is not now generally used for marketing purposes, but is still used internally, and said to reflect the degree of changes to the core of the operating system. However, for application compatibility reasons, Microsoft kept the major version number as 6 in releases following Vista, but changed it later to 10 in Windows 10. The build number is an internal identifier used by Microsoft's developers and beta testers.
Starting with Windows 8.1, Microsoft changed the Version API Helper functions' behavior. If an application is not manifested for Windows 8.1 or later, the API will always return version 6.2, which is the version number of Windows 8. This is because the manifest feature was introduced with Windows 8.1, to replace GetVersion and related functions.
Supported platforms
32-bit platforms
In order to prevent Intel x86-specific code from slipping into the operating system, due to developers being used to developing on x86 chips, Windows NT 3.1 was initially developed using non-x86 development systems and then ported to the x86 architecture. This work was initially based on the Intel i860-based Dazzle system and, later, the MIPS R4000-based Jazz platform. Both systems were designed internally at Microsoft.
Windows NT 3.1 was released for Intel x86 PC compatible and PC-98 platforms, and for DEC Alpha and ARC-compliant MIPS platforms. Windows NT 3.51 added support for the PowerPC processor in 1995, specifically PReP-compliant systems such as the IBM ThinkPad Power Series laptops and Motorola PowerStack series; but despite meetings between Michael Spindler and Bill Gates, not on the Power Macintosh as the PReP compliant Power Macintosh project failed to ship.
Intergraph Corporation ported Windows NT to its Clipper architecture and later announced an intention to port Windows NT 3.51 to Sun Microsystems' SPARC architecture, in conjunction with the company's planned introduction of UltraSPARC models in 1995, but neither version was sold to the public as a retail product.
Only two of the Windows NT 4.0 variants (IA-32 and Alpha) have a full set of service packs available. All of the other ports done by third parties (Motorola, Intergraph, etc.) have few, if any, publicly available updates.
Windows NT 4.0 was the last major release to support Alpha, MIPS, or PowerPC, though development of Windows 2000 for Alpha continued until August 1999, when Compaq stopped support for Windows NT on that architecture; and then three days later Microsoft also canceled their AlphaNT program, even though the Alpha NT 5 (Windows 2000) release had reached RC1 status.
On January 5, 2011, Microsoft announced that the next major version of the Windows NT family will include support for the ARM architecture. Microsoft demonstrated a preliminary version of Windows (version 6.2.7867) running on an ARM-based computer at the 2011 Consumer Electronics Show. This eventually led to the commercial release of the Windows 8-derived Windows RT on October 26, 2012, and the use of Windows NT, rather than Windows CE, in Windows Phone 8.
The original Xbox and Xbox 360 run a custom operating system based upon a heavily modified version of Windows 2000, an approach that Microsoft engineer Don Box called "fork and run". It exports APIs similar to those found in Microsoft Windows, such as Direct3D. The Xbox One and Xbox Series X/S consoles use a stripped-down version of the Windows operating system.
Windows 11 is the first non-server version of Windows NT that does not support 32-bit platforms.
64-bit platforms
The 64-bit versions of Windows NT were originally intended to run on Itanium and DEC Alpha; the latter was used internally at Microsoft during early development of 64-bit Windows. This continued for some time after Microsoft publicly announced that it was cancelling plans to ship 64-bit Windows for Alpha. Because of this, Alpha versions of Windows NT are 32-bit only.
While Windows 2000 only supports Intel IA-32 (32-bit), Windows XP, Server 2003, Server 2008 and Server 2008 R2 each have one edition dedicated to Itanium-based systems. In comparison with Itanium, Microsoft adopted x64 on a greater scale: every version of Windows since Windows XP (which has a dedicated x64 edition) has x64 editions.
The first version of Windows NT to support ARM64 devices with Qualcomm processors was Windows 10, version 1709. This is a full version of Windows, rather than the cut-down Windows RT.
Hardware requirements
The minimum hardware specification required to run each release of the professional workstation version of Windows NT has been fairly slow-moving until the 6.0 (Vista) release, which requires a minimum of 15 GB of free disk space, a tenfold increase in free disk space alone over the previous version, and the 2021 10.0 (11) release which excludes most systems built before 2018.
| Technology | Operating Systems | null |
21291678 | https://en.wikipedia.org/wiki/Windows%203.1 | Windows 3.1 | Windows 3.1 is a major release of Microsoft Windows. It was released to manufacturing on April 6, 1992, as a successor to Windows 3.0. Like its predecessors, the Windows 3.1 series run as a shell on top of MS-DOS; it was the last Windows 16-bit operating environment as all future versions of Windows had moved to 32-bit.
Windows 3.1 introduced the TrueType font system as a competitor to Adobe Type Manager. Its multimedia was also expanded, and screensavers were introduced, alongside new software such as Windows Media Player and Sound Recorder. File Manager and Control Panel received tweaks, while Windows 3.1 also saw the introduction of the Windows Registry and add-ons, and it could utilize more memory than its predecessors.
Microsoft also released special versions of Windows 3.1 throughout 1992 and 1993; in Europe and Japan, Windows 3.1 was introduced with more language support, while Tandy Video Information System received a special version, called Modular Windows. In November 1993, Windows 3.11 was released as a minor update, while Windows 3.2 was released as a Simplified Chinese version of Windows 3.1. Microsoft also introduced Windows for Workgroups, the first version of Windows to allow integrated networking. Mostly oriented towards businesses, it received network improvements and it allowed users to share files, use print servers, and chat online, while it also introduced peer-to-peer networking.
The series is considered to be an improvement on its predecessors. It was praised for its reinvigoration of the user interface and technical design. Windows 3.1 sold over three million copies during the first three months of its release, although its counterpart Windows for Workgroups was noted as a "business disappointment" due to its small amount of sold copies. It was succeeded by Windows 95, and Microsoft ended the support for Windows 3.1 series on December 31, 2001, except for the embedded version, which was retired in 2008.
Development history
Windows 3.0, the predecessor of 3.1, was released in 1990, and is considered to be the first version of Windows to receive critical acclaim. Windows 3.0 received around 10 million sales before the release of Windows 3.1 on April 6, 1992. Microsoft began a television advertising campaign for the first time on March 1, 1992. The advertisements, developed by Ogilvy & Mather, were designed to introduce a broader audience to Windows.
Windows 3.1 was codenamed Janus. Like its predecessors, the operating environment runs as a shell on top of MS-DOS, although it does not include the MS-DOS Executive shell. After the introduction of Windows 1.0, Microsoft had worked on gaining support from companies to expand its operating environment on different types of PCs. Tandy Corporation was open to shipping Tandy Sensation PCs with the Windows 3.1 operating environment. IBM and its PCs were also provided with Windows 3.1.
Release versions and features
Windows 3.1
Further enhancements were introduced in Windows 3.1. The TrueType font system was introduced to provide scalable fonts to Windows applications, without having to resort on using third-party technology such as Adobe Type Manager (ATM). Windows 3.1 introduced Arial, Courier New, and Times New Roman fonts, in regular, bold, italic, and bold-italic versions, which could be scaled to any size and rotated, depending on the application.
To improve user interaction, Microsoft initiated warning and event sounds, and introduced computer command shortcuts for copy, cut, and paste. Windows 3.1 is also noted for its improvement of multimedia; screensavers, Windows Media Player, and Sound Recorder were introduced into the operating environment. These features were already present on the Windows 3.0 with Multimedia Extensions version, although they were only available to users with newly bought PCs. The Media Player could play MIDI music files and AVI video files, while the Sound Recorder could play, record, and edit sound files that were affiliated with the WAV format. Minesweeper was officially introduced in Windows 3.1 as a replacement for Reversi, alongside Solitaire. MS-DOS programs were previously not able to be controlled with a mouse; this ended up being introduced in Windows 3.1. Object Linking and Embedding (OLE) was added to allow drag-and-drop embedding of images and formatted text between Windows programs. SVGA color support was also introduced in this version.
File Manager had also received tweaks; split view-mode was introduced, users were now able to browse files without having to open separate windows, while files were able to be dragged and dropped to other locations on the system. An option for quick formatting was introduced to format floppy disks and copy its files without having to quit Windows. File Manager is an MDI application that is used for moving, deleting, and managing files on the system. Microsoft also built Microsoft Bob, a utility that would act as a search assistant, on Windows 3.1, only for it to be released on Windows 95 in 1995. The introduction of Windows Registry, a centralized database that could store configuration information and settings for various operating systems components and applications, also occurred in this version. The Control Panel also received changes; its items were now hard-coded, and additional items could be added by placing additional .cpl files. Similarly, the Calendar uses the .cal extension. Printer management tasks were moved over to Control Panel and Print Manager. Several printer drivers were improved in Windows 3.1, making the Print Manager more efficient to use. Windows 3.1 also includes troubleshooting and diagnostic tools such as the Dr. Watson utility which saves information about application errors, and Microsoft Diagnostics.
Windows 3.1 also includes add-ons; Video for Windows was introduced in November 1992 as a reaction to Apple's QuickTime technology. At the price of $200, the software included editing and encoding programs. It was later built into Windows 95. Microsoft also published Windows for Pen Computing, a pen computing interface which was created in response to PenPoint OS by GO Corporation. The operating environment was also given limited compatibility with the then-new 32-bit Windows API, by introducing Win32s, an enabling technology. Microsoft also provided WinG, an application program interface, to entice developers to move from DOS to Windows. It also provided a device-independent interface to graphics and printer hardware, and allowed programs to have both read and write capabilities to the WinGDC.
Unlike all previous versions, Windows 3.1 could not run in real mode and it insisted on the use of 80286 processors or above. Because of this, the maximum memory available was increased. While Windows 3.0 was limited to 16 MB maximum memory, Windows 3.1 could access a theoretical 4 GB in the 386 enhanced mode. The actual practical ceiling is 256 MB. Like its predecessors, it runs as a 16-bit system; Windows 3.1 is also the last Windows to run in 16-bit mode. It is also the first Windows to be distributed on a CD-ROM. The setup interface was simplified; express mode was introduced to automatically set up Windows. Windows 3.1 also includes an online tutorial applet for users regarding the use of the Windows 3.1 user interface. In addition it supported the Advanced Power Management standard.
Windows 3.1 for Central and Eastern Europe
A special version named "Windows 3.1 for Central and Eastern Europe" introduced eleven languages to Windows 3.1. It also provided support for the Cyrillic script. To use Czech, Hungarian, and Polish terminologies this version was required, while to use Russian terminologies a Russian version of Windows 3.1 was needed. Similarly, Microsoft also released Windows 3.1J with support for Japanese, which shipped 1.46 million copies in its first year on the market (1993) in Japan.
Modular Windows
Modular Windows was built for real-time consumer electronics, and was designed to be controlled via television. It was a special version of Windows 3.1, which was designed to run on Tandy Video Information System; it allowed users to run multimedia software without having to buy a personal computer. It also contained a software development kit (SDK) for programmers to write applications that would run on devices that have Modular Windows. The SDK was sold for $99. Modular Windows was discontinued in 1994.
Windows 3.11
Released on November 8, 1993, Windows 3.11 was introduced with fixes for network problems which were present on Windows 3.1. As a minor update, new features were not present in this version. It also did not run on IBM's OS/2 for Windows. Windows 3.11 allowed users to connect to each other as peers to share the resources of their computers. Microsoft replaced all retail and OEM versions of Windows 3.1 with Windows 3.11 and provided a free upgrade to anyone who owned Windows 3.1.
Windows 3.2
An updated Simplified Chinese version of Windows 3.1 was released in November 1993, as Windows 3.2. The update was limited to this language version, as it only fixed issues related to the complex input system for the Simplified Chinese language. A font editor is present in Windows 3.2; it is used to add new Chinese characters to the already-existing fonts.
Windows for Workgroups
Windows for Workgroups served as an update to Windows 3.1, and it was the first version of Windows that was suitable for integrated networking. Initially developed as an add-on for Windows 3.0, it was later released in 1992. It introduced drivers and protocols for peer-to-peer networking. Windows for Workgroups was mostly oriented towards businesses.
Windows for Workgroups 3.1
The first version of Windows for Workgroups, 3.1, was released on October 27, 1992. Codenamed Winball and Sparta, it allows users to share files, use print servers, and chat online; files could be accessed from other machines that run either Windows or DOS. The Microsoft Hearts card game was also added, while Object Linking and Embedding, which was implemented in Windows 3.1, was also included in the Windows for Workgroups version. The Workgroups version also introduced the Microsoft Mail program, which allowed users to receive and send email, and Microsoft Schedule+, a time management app.
Windows for Workgroups could also be accessed from an OS/2 client that uses the Server Message Block (SMB), a protocol used for sharing files and printers over local networks. It introduced support for the NetBEUI protocol. The price sat at $69 for Windows 3.1 users.
Windows for Workgroups 3.11
The other version, Windows for Workgroups 3.11, was released on November 8, 1993. It was codenamed Snowball, and it introduced support for 32-bit file access, drive sharing, and group calendaring. It also has built-in fax capabilities.
It received network improvements; a Winsock package was released for Windows for Workgroups, although it was later replaced by a 32-bit stack add-on package (codenamed Wolverine) that provided TCP/IP support in Windows for Workgroups 3.11. Its connectivity with NetWare networks was increased, while it also introduced support for Open Data-Link Interface cards and Internetwork Packet Exchange drivers. Remote access service was introduced as a product for users to remotely access Windows NT and its Advanced Server networks.
It runs in 80386 enhanced mode, and it supports the use of network redirectors. It was sold in two versions; the complete package cost while the "Workgroup Add-on for Windows" cost .
System requirements
The official system requirements for Windows 3.1 and subsequent versions include the following:
To use a printer or to run Windows on a network, additional 2.5 MB of free space will be needed on the hard drive. The amount of RAM is dependent on software that runs on the PC; if the user is on the network and if the network requires a lot of memory, more RAM will be needed. Windows 3.1 includes more drivers for printers than its predecessor. It is also possible to connect to a network using Windows 3.1 via Hayes, Multi-Tech, or Trail Blazer modems.
Reception
Windows 3.1 is considered to be more stable and multimedia-friendly in comparison with its predecessor, while its user interface was reinvigorated. It has been shown as an improvement, and it possesses more features in comparison with its rival IBM OS/2 2.0, which launched a month earlier than Windows 3.1. InfoWorld rated the operating environment a "very good" value. Windows for Workgroups received lukewarm reception; it has been praised for its technical design, but it has been also noted as a "business disappointment" due to its small amount of sold copies.
Regarding the marketplace, Windows 3.1 had received an enthusiastic reception; its retail price sat at $149, and over three million copies of Windows 3.1 were sold in the first three months. The year of Windows 3.1's release was successful for Microsoft, which was named the "Most Innovative Company Operating in the U.S." by Fortune magazine, while Windows became the most widely used GUI-based operating environment.
Microsoft ended its support for Windows 3.1 and Windows for Workgroups on December 31, 2001, although the embedded version of Windows for Workgroups 3.11 was retired on November 1, 2008. The operating environment was superseded by Windows NT 3.1, which was released in 1993, and Windows 95 in 1995.
DR-DOS compatibility
The installer of the beta release used code that checked whether it was running on Microsoft-licensed DOS or another DOS operating system, such as DR-DOS. It was known as AARD code, and Microsoft disabled it before the final release of Windows 3.1, though without removing it altogether. Digital Research, who owned DR-DOS, released a patch within weeks to allow the installer to continue. Memos that were released during the United States v. Microsoft Corp. antitrust case in 1999 revealed that Microsoft specifically focused it on DR-DOS. When Caldera bought DR-DOS from Novell, they brought a lawsuit against Microsoft over the AARD code, which was later settled with Microsoft paying $280 million.
Legacy
Windows 3.1 found a niche market as an embedded operating system after becoming obsolete in the PC world. By 2008, both Virgin Atlantic and Qantas employed it for some of the onboard entertainment systems on long-distance jets. It also sees continued use as an embedded OS in retail cash tills. On July 14, 2013, Linux kernel version 3.11 was officially named "Linux for Workgroups" as a tongue-in-cheek reference to Windows for Workgroups 3.11.
In November 2015, the failure of a Windows 3.1 system in Orly Airport in Paris, which was responsible for communicating visual range information in foggy weather to pilots, made operations temporarily cease. Whether the failure was hardware- or software-based is not specified, though the highlighting of the operating system suggests a software failure. In 2016, the Internet Archive organization released Windows 3.1 as an emulated environment in a web browser.
In January 2024, German state-owned national railway company Deutsche Bahn posted a job listing for a system administrator with "knowledge of legacy operating systems". The main responsibilities listed in the post were maintenance of the old system and driver updates. The need for the continued use of Windows 3.11 could apparently be traced back to Siemens' SIBAS (Siemens Bahn Automatisierungs System) automation system used to control trains. The job post was retracted due to "unfortunate wording".
| Technology | Operating Systems | null |
21291954 | https://en.wikipedia.org/wiki/MS-DOS | MS-DOS | MS-DOS ( ; acronym for Microsoft Disk Operating System, also known as Microsoft DOS) is an operating system for x86-based personal computers mostly developed by Microsoft. Collectively, MS-DOS, its rebranding as IBM PC DOS, and a few operating systems attempting to be compatible with MS-DOS, are sometimes referred to as "DOS" (which is also the generic acronym for disk operating system). MS-DOS was the main operating system for IBM PC compatibles during the 1980s, from which point it was gradually superseded by operating systems offering a graphical user interface (GUI), in various generations of the graphical Microsoft Windows operating system.
IBM licensed and re-released it in 1981 as PC DOS 1.0 for use in its PCs. Although MS-DOS and PC DOS were initially developed in parallel by Microsoft and IBM, the two products diverged after twelve years, in 1993, with recognizable differences in compatibility, syntax and capabilities. Beginning in 1988 with DR-DOS, several competing products were released for the x86 platform.
Initially, MS-DOS was targeted at Intel 8086 processors running on computer hardware using floppy disks to store and access not only the operating system, but application software and user data as well. Progressive version releases delivered support for other mass storage media in ever greater sizes and formats, along with added feature support for newer processors and rapidly evolving computer architectures. Ultimately, it was the key product in Microsoft's development from a programming language company to a diverse software development firm, providing the company with essential revenue and marketing resources. It was also the underlying basic operating system on which early versions of Windows ran as a GUI. MS-DOS went through eight versions, until development ceased in 2000; version 6.22 from 1994 was the final standalone version, with versions 7 and 8 serving mostly in the background for loading Windows 9x.
History
MS-DOS was a renamed form of 86-DOSowned by Seattle Computer Products, written by Tim Paterson. Development of 86-DOS took only six weeks, as it was basically a clone of Digital Research's CP/M (for 8080/Z80 processors), ported to run on 8086 processors and with two notable differences compared to CP/M: an improved disk sector buffering logic, and the introduction of FAT12 instead of the CP/M filesystem. This first version was shipped in August 1980. Microsoft, which needed an operating system for the IBM Personal Computer, hired Tim Paterson in May 1981 and bought 86-DOS 1.10 for in July of the same year. Microsoft kept the version number, but renamed it MS-DOS. They also licensed MS-DOS 1.10/1.14 to IBM, which, in August 1981, offered it as PC DOS 1.0 as one of three operating systems for the IBM 5150 or the IBM PC.
Within a year, Microsoft licensed MS-DOS to over 70 other companies. It was designed to be an OS that could run on any 8086-family computer. Each computer would have its own distinct hardware and its own version of MS-DOS, similar to the situation that existed for CP/M, and with MS-DOS emulating the same solution as CP/M to adapt for different hardware platforms. To this end, MS-DOS was designed with a modular structure with internal device drivers (the DOS BIOS), minimally for primary disk drives and the console, integrated with the kernel and loaded by the boot loader, and installable device drivers for other devices loaded and integrated at boot time. The OEM would use a development kit provided by Microsoft to build a version of MS-DOS with their basic I/O drivers and a standard Microsoft kernel, which they would typically supply on disk to end users along with the hardware. Thus, there were many different versions of "MS-DOS" for different hardware, and there is a major distinction between an IBM-compatible (or ISA) machine and an MS-DOS [compatible] machine. Some machines, like the Tandy 2000, were MS-DOS compatible but not IBM-compatible, so they could run software written exclusively for MS-DOS without dependence on the peripheral hardware of the IBM PC architecture.
This design would have worked well for compatibility, if application programs had only used MS-DOS services to perform device I/O. Indeed, the same design philosophy is embodied in Windows NT (see Hardware Abstraction Layer). However, in MS-DOS' early days, the greater speed attainable by programs through direct control of hardware was of particular importance, especially for games, which often pushed the limits of their contemporary hardware. Very soon an IBM-compatible architecture became the goal, and before long all 8086-family computers closely emulated IBM's hardware, and only a single version of MS-DOS for a fixed hardware platform was needed for the market. This version is the version of MS-DOS that is discussed here, as the dozens of other OEM versions of "MS-DOS" were only relevant to the systems they were designed for, and in any case were very similar in function and capability to some standard version for the IBM PC—often the same-numbered version, but not always, since some OEMs used their own proprietary version numbering schemes (e.g. labeling later releases of MS-DOS 1.x as 2.0 or vice versa)—with a few notable exceptions.
Microsoft omitted multi-user support from MS-DOS because Microsoft's Unix-based operating system, Xenix, was fully multi-user. The company planned, over time, to improve MS-DOS so it would be almost indistinguishable from single-user Xenix, or XEDOS, which would also run on the Motorola 68000, Zilog Z8000, and the LSI-11; they would be upwardly compatible with Xenix, which Byte in 1983 described as "the multi-user MS-DOS of the future". Microsoft advertised MS-DOS and Xenix together, listing the shared features of its "single-user OS" and "the multi-user, multi-tasking, UNIX-derived operating system", and promising easy porting between them. After the breakup of the Bell System, however, AT&T Computer Systems started selling UNIX System V. Believing that it could not compete with AT&T in the Unix market, Microsoft abandoned Xenix, and in 1987 transferred ownership of Xenix to the Santa Cruz Operation (SCO).
On March 25, 2014, Microsoft made the code to SCP MS-DOS 1.25 and a mixture of Altos MS-DOS 2.11 and TeleVideo PC DOS 2.11 available to the public under the Microsoft Research License Agreement, which makes the code source-available, but not open source as defined by Open Source Initiative or Free Software Foundation standards. Microsoft would later re-license the code under the MIT License on September 28, 2018, making these versions free software. Microsoft later released the code for MS-DOS 4.00 on April 25, 2024, under the same license.
As an April Fool's Day joke in 2015, Microsoft Mobile launched a Windows Phone application called MS-DOS Mobile which was presented as a new mobile operating system and worked similar to MS-DOS.
Versions
Microsoft licensed or released versions of MS-DOS under different names like Lifeboat Associates "Software Bus 86" a.k.a. SB-DOS, COMPAQ-DOS, NCR-DOS or Z-DOS before it eventually enforced the MS-DOS name for all versions but the IBM one, which was originally called "IBM Personal Computer DOS", later shortened to IBM PC DOS. (Competitors released compatible DOS systems such as DR-DOS and PTS-DOS that could also run MS-DOS applications.)
In the former Eastern bloc, MS-DOS derivatives named DCP () 3.20 and 3.30 (DCP 1700, DCP 3.3) and WDOS existed in the late 1980s. They were produced by the East German electronics manufacturer VEB Robotron.
The following versions of MS-DOS were released to the public:
MS-DOS 1.x
Version 1.23 (OEM)
Version 1.24 (OEM) – basis for IBM's Personal Computer DOS 1.1
Version 1.25 (OEM) – basis for non-IBM OEM versions of MS-DOS, including SCP MS-DOS 1.25
Compaq-DOS 1.12, a Compaq OEM version of MS-DOS 1.25; Release date: November 1983
TI BOOT V. 1.13, a Texas Instruments OEM version of MS-DOS; Release date: August 1983
Zenith Z-DOS 1.19, a Zenith OEM version of MS-DOS 1.25
Zenith Z-DOS/MS-DOS release 1.01, version 1.25, a Zenith OEM version of MS-DOS; Release date: May 1983
MS-DOS 2.x
Support for IBM's XT 10 MB hard disk drives, support up to 16 MB or 32 MB FAT12-formatted hard disk drives depending on the formatting tool shipped by OEMs, user-installable device drivers, tree-structure filing system, Unix-like inheritable redirectable file handles, non-multitasking child processes an improved Terminate and Stay Resident (TSR) API, environment variables, device driver support, FOR and GOTO loops in batch files, ANSI.SYS.
Version 2.0 (OEM), First version to support double-sided 360 KB 5.25-inch floppy disks; Release date: October 1983
Version 2.02 (OEM, Compaq); Release date: November 1983
Version 2.05 (OEM, international support); Release date: October 1983
Version 2.1 (OEM, IBM only)
Version 2.11 (OEM)
Altos MS-DOS 2.11, an Altos OEM version of MS-DOS 2.11 for the ACT-86C
ITT Corporation ITT-DOS 2.11 Version 2 (MS-DOS 2.11 for the ITT XTRA Personal Computer); Release date: July 1985
Olivetti M19 came with MS-DOS 2.11
Tandy 1000 HX has MS-DOS 2.11 in ROM
TeleVideo PC DOS 2.11, a TeleVideo OEM version of MS-DOS 2.11
Toshiba MS-DOS 2.11 in ROM drive for the model T1000 laptop
Version 2.13 (OEM, Zenith); Release date: July 1984
Version 2.2 (OEM, with Hangeul support)
Version 2.25 (OEM, with Hangeul and Kanji support)
Version 2.3 (used on the Toshiba Pasopia 16)
MS-DOS 3.x
Version 3.0 (OEM) – First version to support 5.25-inch 1.2 MB floppy drives and diskettes, FAT16 partitions up to 32 MB; Release date: April 1985
Version 3.1 (OEM) – Support for Microsoft Networks through an IFS layer, remote file and printer API
Version 3.2 (OEM) – First version to support 3.5-inch 720 KB floppy drives and diskettes and XCOPY.
Version 3.10 (OEM, Multitech); Release date: May 1986
Version 3.20 – First retail release (non-OEM); Release date: July 1986
Version 3.21 (OEM / non-OEM); Release date: May 1987
Version 3.22 (OEM) – (HP 95LX)
Version 3.3 (OEM) – First version to support 3.5-inch 1.44 MB floppy drives and diskettes, extended and logical partitions, directory tree copying with XCOPY, improved support for internationalization (COUNTRY.SYS), networked file flush operations
Version 3.3a (OEM)
Version 3.30; Release date: February 1988
Version 3.30A (OEM, DTK); Release date: July 1987
Version 3.30T (OEM, Tandy); Release date: July 1990
Version 3.31 (Compaq OEM only) – supports FAT16B with partitions larger than 32 MiB; Release date: November 1989
MS-DOS 4.0 / MS-DOS 4.x
MS-DOS 4.0 (multitasking) and MS-DOS 4.1A separate branch of development with additional multitasking features, released between 3.2 and 3.3, and later abandoned. It is unrelated to any later versions, including versions 4.00 and 4.01 listed below
MS-DOS 4.x (IBM-developed) – Includes a graphical/mouse interface. It had many bugs and compatibility issues that plagued all versions of MS-DOS 4.x.
Version 4.00 (OEM) – First version with built-in IBM/Microsoft support for hard disk partitions larger than 32 MB and up to a maximum size of 2 GB, FASTOPEN/FASTSEEK, DOSSHELL, could use EMS for the disk buffers and provided EMS drivers and emulation for 386-compatible processors; Release date: October 1988
Version 4.01 (OEM) – Microsoft-rewritten Version 4.00 released under MS-DOS label but not IBM PC DOS. First version to introduce volume serial number when formatting hard disks and floppy disks (Disk duplication also and when using SYS to make a floppy disk or a hard drive partition bootable); Release date: April 1989
Version 4.01a (OEM)
MS-DOS 5.x
Version 5.0 (Retail) – includes a full-screen text editor. Many of the bugs and compatibility issues from MS-DOS 4.x are resolved. First version to support 3.5-inch 2.88 MB floppy drives and diskettes. The SHARE command was not needed anymore for old DOS 1.x style FCB file API to partitions over 32 MB. First version to get the HIMEM.SYS driver and load portions of the operating system into the upper memory area and high memory area. Supports up to four DOS primary partitions, however FDISK cannot create more than one primary partition. Third-party tools allowed for the creation of up to four primary partitions.
AST Premium Exec DOS 5.0 (OEM)a version for the AST Premium Exec series of notebooks with various extensions, including improved load-high and extended codepage support
Version 5.0a (Retail) – With this release, IBM and Microsoft versions diverge.
Version 5.50 (Windows NTVDM) – All Windows NT 32-bit versions ship with files from DOS 5.0
MS-DOS 6.x
Version 6.0 (Retail) – Online help through QBasic. Disk compression, upper memory optimization and antivirus included.
Version 6.2 – SCANDISK as replacement for CHKDSK. Fix serious bugs in DBLSPACE.
Version 6.21 (Retail) – Stacker-infringing DBLSPACE removed.
Version 6.22 (Retail) – New DRVSPACE compression. Last version of MS-DOS to be sold as an independent product.
MS-DOS 7/8 (as part of Windows 9x)
MS-DOS 7.0 was included in Windows 95's first retail release. It contains support for VFAT long file names when run in a Windows Virtual 8086 box or with an LFN driver such as DOSLFN. JO.SYS is an alternative filename of the IO.SYS kernel file and used as such for "special purposes". JO.SYS allows booting from either CD-ROM drive or hard disk. Last version to recognize only the first 8.4 GB of a hard disk. The VER internal command reports the Windows version 4.00.950, applications through the MS-DOS API would be reported a version number of 7.00.
MS-DOS 7.1 was included in Windows 95's OEM Service Release 2 through Windows 98 Second Edition. It added support for the FAT32 file system and logical block addressing (LBA), and was the last version that could boot to the command line from a hard disk. The VER internal command reports the Windows version 4.00.1111, 4.10.1998, or 4.10.2222 depending on the version of Windows, while applications through the API would report version 7.10.
MS-DOS 8.0 was included in Windows Me, the last version based on MS-DOS. DOS mode was significantly altered in this release. Booting from the hard disk to a command line only was no longer permitted, AUTOEXEC.BAT and CONFIG.SYS files were no longer loaded or parsed before loading the Windows GUI; booting from floppy disk was still permitted to allow for emergency recovery. This version was included (in modified form) in Windows XP up to Windows 8.1 for creating MS-DOS startup disks. The VER internal command reports the Windows version 4.90.3000 or 5.1 when created from newer versions of Windows. Applications requesting the version through the API would report version 8.00.
Microsoft DOS was released through the OEM channel, until Digital Research released DR-DOS 5.0 as a retail upgrade. With PC DOS 5.00.1, the IBM–Microsoft agreement started to end, and IBM entered the retail DOS market with IBM DOS 5.00.1, 5.02, 6.00 and PC DOS 6.1, 6.3, 7, 2000 and 7.1.
Localized versions
Localized versions of MS-DOS existed for different markets. While Western issues of MS-DOS evolved around the same set of tools and drivers just with localized message languages and differing sets of supported codepages and keyboard layouts, some language versions were considerably different from Western issues and were adapted to run on localized PC hardware with additional BIOS services not available in Western PCs, support multiple hardware codepages for displays and printers, support DBCS, alternative input methods and graphics output. Affected issues include Japanese (DOS/V), Korean, Arabic (ADOS 3.3/5.0), Hebrew (HDOS 3.3/5.0), Russian (RDOS 4.01/5.0) as well as some other Eastern European versions of DOS.
Competition
On microcomputers based on the Intel 8086 and 8088 processors, including the IBM PC and clones, the initial competition to the PC DOS/MS-DOS line came from Digital Research, whose CP/M operating system had inspired MS-DOS. In fact, there remains controversy as to whether QDOS was more or less plagiarized from early versions of CP/M code. Digital Research released CP/M-86 a few months after MS-DOS, and it was offered as an alternative to MS-DOS and Microsoft's licensing requirements, but at a higher price. Executable programs for CP/M-86 and MS-DOS were not interchangeable with each other; many applications were sold in both MS-DOS and CP/M-86 versions until MS-DOS became preponderant (later Digital Research operating systems could run both MS-DOS and CP/M-86 software). MS-DOS originally supported the simple .COM, which was modeled after a similar but binary-incompatible format known from CP/M-80. CP/M-86 instead supported a relocatable format using the filename extension .CMD to avoid name conflicts with CP/M-80 and MS-DOS .COM files. MS-DOS version 1.0 added a more advanced relocatable .EXE executable file format.
Most of the machines in the early days of MS-DOS had differing system architectures and there was a certain degree of incompatibility, and subsequently vendor lock-in. Users who began using MS-DOS with their machines were compelled to continue using the version customized for their hardware, or face trying to get all of their proprietary hardware and software to work with the new system.
In the business world, the 808x-based machines that MS-DOS was tied to faced competition from the Unix operating system; the latter ran on many different hardware architectures. Microsoft itself sold a version of Unix for the PC called Xenix.
In the emerging world of home users, a variety of other computers based on various other processors were in serious competition with the IBM PC: the Apple II, Mac, Commodore 64 and others did not use the 808x processor; many 808x machines of different architectures used custom versions of MS-DOS. At first all these machines were in competition. In time the IBM PC hardware configuration became dominant in the 808x market as software written to communicate directly with the PC hardware without using standard operating system calls ran much faster, but on true PC-compatibles only. Non-PC-compatible 808x machines were too small a market to have fast software written for them alone, and the market remained open only for IBM PCs and machines that closely imitated their architecture, all running either a single version of MS-DOS compatible only with PCs, or the equivalent IBM PC DOS. Most clones cost much less than IBM-branded machines of similar performance, and became widely used by home users, while IBM PCs had a large share of the business computer market.
Microsoft and IBM together began what was intended as the follow-on to MS-DOS/PC DOS, called OS/2. When OS/2 was released in 1987, Microsoft began an advertising campaign announcing that "DOS is Dead" and stating that version 4 was the last full release. OS/2 was designed for efficient multi-tasking and offered a number of advanced features that had been designed together with similar look and feel; it was seen as the legitimate heir to the "kludgy" DOS platform.
MS-DOS had grown in spurts, with many significant features being taken or duplicated from Microsoft's other products and operating systems. MS-DOS also grew by incorporating, by direct licensing or feature duplicating, the functionality of tools and utilities developed by independent companies, such as Norton Utilities, PC Tools (Microsoft Anti-Virus), QEMM expanded memory manager, Stacker disk compression, and others.
During the period when Digital Research was competing in the operating system market some computers, like the Amstrad PC1512, were sold with floppy disks for two operating systems (only one of which could be used at a time), MS-DOS and CP/M-86 or a derivative of it. Digital Research produced DOS Plus, which was compatible with MS-DOS 2.11, supported CP/M-86 programs, had additional features including multi-tasking, and could read and write disks in CP/M and MS-DOS format.
While OS/2 was under protracted development, Digital Research released the MS-DOS compatible DR-DOS 5.0, which included features only available as third-party add-ons for MS-DOS. Unwilling to lose any portion of the market, Microsoft responded by announcing the "pending" release of MS-DOS 5.0 in May 1990. This effectively killed most DR DOS sales until the actual release of MS-DOS 5.0 in June 1991. Digital Research brought out DR DOS 6.0, which sold well until the "pre-announcement" of MS-DOS 6.0 again stifled the sales of DR DOS.
Microsoft had been accused of carefully orchestrating leaks about future versions of MS-DOS in an attempt to create what in the industry is called FUD (fear, uncertainty, and doubt) regarding DR DOS. For example, in October 1990, shortly after the release of DR DOS 5.0, and long before the eventual June 1991 release of MS-DOS 5.0, stories on feature enhancements in MS-DOS started to appear in InfoWorld and PC Week. Brad Silverberg, then Vice President of Systems Software at Microsoft and general manager of its Windows and MS-DOS Business Unit, wrote a forceful letter to PC Week (November 5, 1990), denying that Microsoft was engaged in FUD tactics ("to serve our customers better, we decided to be more forthcoming about version 5.0") and denying that Microsoft copied features from DR DOS:
"The feature enhancements of MS-DOS version 5.0 were decided and development was begun long before we heard about DR DOS 5.0. There will be some similar features. With 50 million MS-DOS users, it shouldn't be surprising that DRI has heard some of the same requests from customers that we have." – (Schulman et al. 1994).
The pact between Microsoft and IBM to promote OS/2 began to fall apart in 1990 when Windows 3.0 became a marketplace success. Many of Microsoft's further contributions to OS/2 also went into creating a third GUI replacement for DOS, Windows NT.
IBM, which had already been developing the next version of OS/2, carried on development of the platform without Microsoft and sold it as the alternative to DOS and Windows.
Legal issues
As a response to Digital Research's DR DOS 6.0, which bundled SuperStor disk compression, Microsoft opened negotiations with Stac Electronics, vendor of the most popular DOS disk compression tool, Stacker. In the due diligence process, Stac engineers had shown Microsoft part of the Stacker source code. Stac was unwilling to meet Microsoft's terms for licensing Stacker and withdrew from the negotiations. Microsoft chose to license Vertisoft's DoubleDisk, using it as the core for its DoubleSpace disk compression.
MS-DOS 6.0 and 6.20 were released in 1993, both including the Microsoft DoubleSpace disk compression utility program. Stac successfully sued Microsoft for patent infringement regarding the compression algorithm used in DoubleSpace. This resulted in the 1994 release of MS-DOS 6.21, which had disk compression removed. Shortly afterwards came version 6.22, with a new version of the disk compression system, DriveSpace, which had a different compression algorithm to avoid the infringing code.
Prior to 1995, Microsoft licensed MS-DOS (and Windows) to computer manufacturers under three types of agreement: per-processor (a fee for each system the company sold), per-system (a fee for each system of a particular model), or per-copy (a fee for each copy of MS-DOS installed). The largest manufacturers used the per-processor arrangement, which had the lowest fee. This arrangement made it expensive for the large manufacturers to migrate to any other operating system, such as DR DOS. In 1991, the U.S. government Federal Trade Commission began investigating Microsoft's licensing procedures, resulting in a 1994 settlement agreement limiting Microsoft to per-copy licensing. Digital Research did not gain by this settlement, and years later its successor in interest, Caldera, sued Microsoft for damages in the Caldera v. Microsoft lawsuit. It was believed that the settlement ran in the order of , but was revealed in November 2009 with the release of the Settlement Agreement to be .
Use of undocumented APIs
Microsoft also used a variety of tactics in MS-DOS and several of their applications and development tools that, while operating perfectly when running on genuine MS-DOS (and PC DOS), would break when run on another vendor's implementation of DOS. Notable examples of this practice included:
Microsoft's QuickPascal (released in early 1989) was the first MS product that checked for MS-DOS by modifying the program's Program Segment Prefix using undocumented DOS functions, and then checked whether or not the associated value changed in a fixed position within the DOS data segment (also undocumented). This check also made it into later MS products, including Microsoft QuickC v2.5, Programmer's Workbench and Microsoft C v6.0.
The AARD code, a block of code in the Windows launcher (WIN.COM) and a few other system files of Windows 3.1. It was XOR encrypted, self-modifying, and deliberately obfuscated, using various undocumented DOS structures and functions to determine whether or not Windows really was running on MS-DOS. In the beta versions, it displayed an "error" message if the test for genuine MS-DOS failed, prompting the user to abort or continue, with abort the default. In the final release version, the code still ran, but the message and prompt were disabled by an added flag byte, rendering it (probably) ineffectual.
Windows 3.0 beta release only gave a warning that Windows would not operate properly on a "foreign" OS and actually ran well on DR DOS 6.0.
Interrupt routines called by Windows to inform MS-DOS that Windows is starting/exiting, information that MS-DOS retained in an IN_WINDOWS flag, in spite of the fact that MS-DOS and Windows were supposed to be two separate products.
Windows command-line interface
All versions of Microsoft Windows have had an MS-DOS or MS-DOS-like command-line interface called MS-DOS Prompt which redirected input to MS-DOS and output from MS-DOS to the MS-DOS Prompt, or, in later versions, Command Prompt. This could run many DOS and variously Win32, OS/2 1.x and POSIX command-line utilities in the same command-line session, allowing piping between commands. The user interface, and the icon up to Windows 2000, followed the native MS-DOS interface. The Command Prompt introduced with Windows NT is not actually MS-DOS, but shares some commands with MS-DOS.
Earlier versions of Windows
The 16-bit versions of Windows (up to 3.11) ran as a graphical user interface (GUI) on top of MS-DOS. With Windows 95, 98, and Me, the role of MS-DOS was reduced to a boot loader according to Microsoft, with MS-DOS programs running in a virtual DOS machine within 32-bit Windows, with ability to boot directly into MS-DOS retained as a backward compatibility option for applications that required real mode access to the hardware, which was generally not possible within Windows. The command line accessed the DOS command line (usually COMMAND.COM) through a Windows module (WINOLDAP.MOD).
Windows NT
Windows NT-based operating systems boot to a kernel whose purpose is to load Windows and run the system. One cannot run Win32 applications in the loader system in the manner that OS/2, UNIX or consumer versions of Windows can launch character-mode sessions.
The command session permits running various supported command-line utilities from Win32, MS-DOS, OS/2 1.x and POSIX. The emulators for MS-DOS, OS/2 and POSIX use the host's window in the same way that Win16 applications use the Win32 explorer. Using the host's window allows one to pipe output between emulations.
The MS-DOS emulation takes place through the NTVDM (NT Virtual DOS Machine). This is a modified SoftPC (a former product similar to VirtualPC), running a modified MS-DOS 5 (NTIO.SYS and NTDOS.SYS). The output is handled by the console DLLs, so that the program at the prompt (CMD.EXE, 4NT.EXE, TCC.EXE), can see the output. 64-bit Windows has neither the DOS emulation, nor the DOS commands EDIT, DEBUG and EDLIN that come with 32-bit Windows.
The DOS version returns 5.00 or 5.50, depending on which API function is used to determine it. Utilities from MS-DOS 5.00 run in this emulation without modification. The very early beta programs of NT show MS-DOS 30.00, but programs running in MS-DOS 30.00 would assume that OS/2 was in control.
The OS/2 emulation is handled through OS2SS.EXE and OS2.EXE, and DOSCALLS.DLL. OS2.EXE is a version of the OS/2 shell (CMD.EXE), which passes commands down to the OS2SS.EXE, and input-output to the Windows NT shell. Windows 2000 was the last version of NT to support OS/2. The emulation is OS/2 1.30.
POSIX is emulated through the POSIX shell, but no emulated shell; the commands are handled directly in CMD.EXE.
The Command Prompt is often called the MS-DOS Prompt. In part, this was the official name for it in Windows 9x and early versions of Windows NT (NT 3.5 and earlier), and in part because the SoftPC emulation of DOS redirects output into it. Actually only COMMAND.COM and other 16-bit commands run in an NTVDM with AUTOEXEC.NT and CONFIG.NT initialization determined by _DEFAULT.PIF, optionally permitting the use of Win32 console applications and internal commands with an NTCMDPROMPT directive.
Win32 console applications use CMD.EXE as their command prompt shell. This confusion does not exist under OS/2 because there are separate DOS and OS/2 prompts, and running a DOS program under OS/2 will launch a separate DOS window to run the application.
All versions of Windows for Itanium (no longer sold by Microsoft) and x86-64 architectures no longer include the NTVDM and can therefore no longer natively run DOS or 16-bit Windows applications. There are alternatives such as virtual machine emulators such as Microsoft's own Virtual PC, as well as VMware, DOSBox etc., unofficial compatibility layers such as NTVDMx64, OTVDM (WineVDM), Win3mu and others.
End-of-life
The introduction of Windows 3.0 in 1990, with an easy-to-use graphical user interface, marked the beginning of the end for the command-line driven MS-DOS. With the release of Windows 95 (and continuing in the Windows 9x product line through to Windows Me), an integrated version of MS-DOS was used for bootstrapping, troubleshooting, and backwards-compatibility with old DOS software, particularly games, and no longer released as a standalone product. In Windows 95, the DOS, called MS-DOS 7, can be booted separately, without the Windows GUI; this capability was retained through Windows 98 Second Edition. Windows Me removed the capability to boot its underlying MS-DOS 8.0 alone from a hard disk, but retained the ability to make a DOS boot floppy disk (called an "Emergency Boot Disk") and can be hacked to restore full access to the underlying DOS. On December 31, 2001, Microsoft declared all versions of MS-DOS 6.22 and older obsolete and stopped providing support and updates for the system. As MS-DOS 7.0 was a part of Windows 95, support for it also ended when Windows 95 extended support ended on December 31, 2001. As MS-DOS 7.10 and MS-DOS 8.0 were part of Windows 98 and Windows ME, respectively, support ended when Windows 98 and ME extended support ended on July 11, 2006, thus ending support and updates of MS-DOS from Microsoft.
In contrast to the Windows 9x series, the Windows NT-derived 32-bit operating systems (Windows NT, 2000, XP and newer), developed alongside the 9x series, do not contain MS-DOS compatibility as a core component of the operating system nor do they rely on it for bootstrapping, as NT was not with the level of support for legacy MS-DOS and Win16 apps that Windows 9x was, but does provide limited DOS emulation called NTVDM (NT Virtual DOS Machine) to run DOS applications and provide DOS-like command prompt windows. 64-bit versions of Windows NT prior to Windows 11 (and Windows Server 2008 R2 by extension) do not provide DOS emulation and cannot run DOS applications natively. Windows XP contains a copy of the Windows Me boot disk, stripped down to bootstrap only. This is accessible only by formatting a floppy as an "MS-DOS startup disk". Files like the driver for the CD-ROM support were deleted from the Windows Me bootdisk and the startup files (AUTOEXEC.BAT and CONFIG.SYS) no longer had content. This modified disk was the base for creating the MS-DOS image for Windows XP. Some of the deleted files can be recovered with an undelete tool. When booting up an MS-DOS startup disk made with Windows XP's format tool, the version number and the VER internal command reports as "Windows Millennium" and "5.1", respectively, and not as "MS-DOS 8.0" (which was used as the base for Windows Me but never released as a stand-alone product), though the API still says Version 8.0.
The creation of the MS-DOS startup disk was then carried over to later versions of Windows, with the majority of its contents remaining unchanged from its introduction in Windows XP. When creating a DOS startup disk on Windows Vista, the files on the startup disk are dated April 18, 2005, but are otherwise unchanged, including the string "MS-DOS Version 8 Copyright 1981–1999 Microsoft Corp" inside COMMAND.COM. Windows 7, 8, and 8.1 can also create a MS-DOS startup disk. Starting with Windows 10, the ability to create a MS-DOS startup disk has been removed, and so either a virtual machine running MS-DOS or an older version (in a virtual machine or dual boot) must be used to format a floppy disk, or an image must be obtained from an external source. Other solutions include using DOS compatible alternatives, such as FreeDOS or even copying the required files and boot sector themselves. The last remaining components related to MS-DOS was the NTVDM component, which was removed entirely in Windows starting with Windows 11 as the operating system dropped support for 32-bit processors in favor of being solely offered in 64-bit versions only. This effectively ended any association of MS-DOS within Microsoft Windows after 36 years.
MS-DOS 6.22 was the last standalone version produced by Microsoft for Intel 8088, Intel 8086, and Intel 80286 processors, which remains available for download via their MSDN, volume license, and OEM license partner websites, for customers with valid login credentials. MS-DOS is still used in embedded x86 systems due to its simple architecture and minimal memory and processor requirements, though some current products have switched to the still-maintained open-source alternative FreeDOS.
In 2018, Microsoft released the source code for MS-DOS 1.25 and 2.0 on GitHub, with the source code for MS-DOS 4.00 being released in the same repository six years later. The purpose of this, according to Microsoft, is mainly for education and experimentation with historic operating systems and for new programmers to gain an understanding of how low-level software works, both historic and current. According to program manager Rich Turner, the other versions could not be open-sourced due to third-party licensing restrictions.
Due to the historical nature of the software, Microsoft will not accept any pull requests to the code. Users, however, are allowed and fully encouraged to fork the repository containing the MS-DOS source code and make their own modifications, and do whatever they like with it.
Legacy compatibility
From 1983 onwards, various companies worked on graphical user interfaces (GUIs) capable of running on PC hardware. However, this required duplicated effort and did not provide much consistency in interface design (even between products from the same company).
Later, in 1985, Microsoft Windows 1.0 was released as Microsoft's first attempt at providing a consistent user interface (for applications). The early versions of Windows ran on top of MS-DOS. At first Windows met with little success, but this was also true for most other companies' efforts as well, for example GEM. After version 3.0, Windows gained market acceptance.
Windows 9x used MS-DOS to boot the Windows kernel in protected mode. Basic features related to the file system, such as long file names, were only available to DOS applications when running through Windows. Windows NT runs independently of DOS but includes NTVDM, a component for simulating a DOS environment for legacy applications. It was not included with Windows 11 as the operating system is exclusively offered in 64-bit architectures such as x86-64.
Related systems
MS-DOS compatible systems include:
IBM PC DOS
DR-DOS, Novell DOS, OpenDOS
FreeDOS
PTS-DOS
ROM-DOS
Microsoft made IBM PC DOS for IBM. It and MS-DOS were identical products that eventually diverged starting with MS-DOS version 6.0. Digital Research did not follow Microsoft's version numbering scheme. For example, MS-DOS 4, released in July 1988, was followed by DR DOS 5.0 in May 1990. MS-DOS 5.0 came in April 1991, and DR DOS 6.0 was released the following June.
These products are collectively referred to as "DOS", even though "Disk Operating System" is a generic term used on other systems unrelated to the x86 and IBM PC. "MS-DOS" can also be a generic reference to DOS on IBM PC compatible computers.
Microsoft's control of the Windows platform, and their programming practices which intentionally made Windows appear as if it ran poorly on competing versions of DOS, crippled the ability of other DOS makers to continue to compete with MS-DOS. Digital Research had to release interim releases to circumvent Windows limitations inserted artificially, designed specifically to provide Microsoft with a competitive advantage.
| Technology | Operating systems | null |
2005663 | https://en.wikipedia.org/wiki/Digitaria | Digitaria | Digitaria is a genus of plants in the grass family native to tropical and warm temperate regions but can occur in tropical, subtropical, and cooler temperate regions as well. Common names include crabgrass, finger-grass, and fonio. They are slender monocotyledonous annual and perennial lawn, pasture, and forage plants; some are often considered lawn pests. Digitus is the Latin word for "finger", and they are distinguished by the long, finger-like inflorescences they produce.
Uses
The seeds are edible, most notably those of fonio (Digitaria exilis and Digitaria iburua), Digitaria sanguinalis, as well as Digitaria compacta. They can be toasted, ground into a flour, made into porridge or fermented to make beer. Fonio has been widely used as a staple crop in parts of Africa. It also has decent nutrient qualities as a forage for cattle.
Lawns
The prevalent species of Digitaria in North America are large crabgrass (D. sanguinalis), sometimes known as hairy crabgrass; and smooth crabgrass (D. ischaemum). These species often become problem weeds in lawns and gardens, growing especially well in thin lawns that are watered lightly, under-fertilized, and poorly drained. They are annual plants, and one plant is capable of producing 150,000 seeds per season. The seeds germinate in the late spring and early summer and outcompete the domesticated lawn grasses, expanding outward in a circle up to in diameter. In the autumn when the plants die, they leave large voids in the lawn. The voids then become prime areas for the crabgrass seeds to germinate the following season.
Biological control is preferable over herbicide use on lawns, as crabgrass emergence is not the cause of poor lawn health but a symptom, and it will return annually if the lawn is not restored with fertilization and proper watering. Crabgrass is quickly outcompeted by healthy lawn grass because, as an annual plant, crabgrass dies off in autumn and needs open conditions for its germination the following spring.
Selected species
Digitaria ammophila (Benth.) Hughes – silky umbrella grass
Digitaria bicornis (Lam.) Roemer & J.A.Schultes ex Loud. – Asian crabgrass
Digitaria brownii (Roem. & Schult.) Hughes – cotton panic grass
Digitaria californica (Benth.) Henrard – Arizona cottontop
Digitaria ciliaris (Retz.) Koeler – summer grass, southern crabgrass
Digitaria cognata (Schult.) Pilg. – fall witchgrass, Carolina crabgrass, mountain hairgrass
Digitaria compacta (Roth ex Roem. & Schult.) Veldkamp
Digitaria cruciata (Nees) A. Camus
Digitaria ctenantha (F.Muell.) Hughes – comb finger grass
Digitaria didactyla Willd. – Queensland blue couch
Digitaria divaricatissima var. macractinia (Benth.) Heather L.Stewart & N.G.Walsh
Digitaria eriantha Steud. – pangolagrass, Smuts finger grass, woolly finger grass
Digitaria exilis (Kippist) Stapf – white fonio
Digitaria filiformis (L.) Koeler – slender crabgrass
Digitaria gracillima (Scribn.) Fernald – longleaf crabgrass
Digitaria horizontalis Willd. – Jamaican crabgrass
Digitaria iburua Stapf – Black fonio
Digitaria insularis (L.) Fedde – sourgrass
Digitaria ischaemum (Schreb.) Schreb. ex Muhl. – smooth crabgrass, small crabgrass, smooth finger grass
Digitaria longiflora (Retz.) Pers. – Indian crabgrass
Digitaria milanjiana (Rendle) Stapf – Madagascar crabgrass
Digitaria nuda Schumacher – naked crabgrass
Digitaria pauciflora Hitchc. – twospike crabgrass
Digitaria platycarpha (Trin.) Stapf
Digitaria petelotii Henrard
Digitaria radicosa (C.Presl) Miq. – trailing crabgrass
Digitaria sanguinalis (L.) Scop. – hairy crabgrass, large crabgrass
Digitaria seriata Stapf – Sandveld fingergrass
Digitaria setigera Roth ex Roemer & J.A.Schultes – East Indian crabgrass
Digitaria serotina (Walt.) Michx. – dwarf crabgrass
Digitaria stenotaphrodes (Nees ex Steud.) Stapf
Digitaria texana A.S.Hitchc. – Texas crabgrass
Digitaria velutina (Forssk.) P.Beauv – velvet crabgrass
Digitaria villosa (Walt.) Pers. – shaggy crabgrass
Digitaria violascens Link – violet crabgrass
| Biology and health sciences | Poales | Plants |
2007102 | https://en.wikipedia.org/wiki/Mesonychia | Mesonychia | Mesonychia ("middle claws") is an extinct taxon of small- to large-sized carnivorous ungulates related to artiodactyls. Mesonychians first appeared in the early Paleocene, went into a sharp decline at the end of the Eocene, and died out entirely when the last genus, Mongolestes, became extinct in the early Oligocene. In Asia, the record of their history suggests they grew gradually larger and more predatory over time, then shifted to scavenging and bone-crushing lifestyles before the group became extinct.
Mesonychians probably originated in China, where the most primitive mesonychian, Yantanglestes, is known from the early Paleocene. They were also most diverse in Asia, where they occur in all major Paleocene faunas. Since other predators, such as creodonts and Carnivora, were either rare or absent in these animal communities, mesonychians most likely dominated the large predator niche in the Paleocene of eastern Asia.
One genus, Dissacus, had successfully spread to Europe and North America by the early Paleocene. Dissacus was a jackal-sized predator that has been found all over the Northern Hemisphere, but species of a closely related or identical genus, Ankalagon, from the early to middle Paleocene of New Mexico, were far larger, growing to the size of a bear. A later genus, Pachyaena, entered North America by the earliest Eocene, where it evolved into species that were at least as large. Mesonychians in North America were by far the largest predatory mammals during the early Paleocene to middle Eocene.
Characteristics
Mesonychians have often been reconstructed as resembling wolves albeit superficially, but they would have appeared very different in life. With a short lower spine stiffened by revolute joints, they would have run with stiff backs like modern ungulates rather than bounding or loping with flexible spines like modern carnivorans. While later mesonychians evolved a suite of limb adaptations for running similar to those in both wolves and deer, their legs remained comparatively thick. They would have resembled no group of living animals. Early mesonychians probably walked on the flats of their feet (plantigrade), while later ones walked on their toes (digitigrade). These later mesonychians had hooves, one on each toe, with four toes on each foot. The foot was compressed for efficient running with the axis between the third and fourth toes (paraxonic); it would have looked something like a hoofed paw.
Mesonychians varied in size; some species were as small as a fox, others as large as a horse. Some members of the group are known only from skulls and jaws, or have fragmentary postcranial remains. But where skeletons are known, they indicate that mesonychians had large heads with strong jaw muscles, relatively long necks, and robust bodies with robust limbs that could run effectively but not rotate the hand or reach out to the side. An unrelated early group of mammalian predators, the creodonts, also had unusually large heads and limbs that traded flexibility for efficiency in running; large head size may be connected to inability to use the feet and claws to help catch and process food, as many modern carnivorans do. Some mesonychians are reconstructed as predatory (comparable to canids), others as scavengers or carnivore-scavengers with bone-crushing adaptations to their teeth (comparable to the large hyenas), and some as omnivorous (comparable to pigs, humans, or black bears). They may not have included hypercarnivores (comparable to felids); their teeth were not as effective at cutting meat as later groups of large mammalian predators. In some localities, multiple species or genera coexisted in different ecological niches. There is evidence to suggest that some genera were sexually dimorphic. Some genera may need revision to clarify the actual number of species or remove ambiguity about genera (such as Dissacus and Ankalagon).
These "wolves on hooves" were probably one of the more important predator groups in the late Paleocene and Eocene ecosystems of Europe (which was an archipelago at the time), Asia (which was an island continent), and North America. Mesonychian dentition consisted of molars modified to generate vertical shear, thin blade-like lower molars, and carnassial notches, but no true carnassials. The molars were laterally compressed and often blunt, and were probably used for shearing meat or crushing bones. The largest species are considered to have been scavengers. Many species are suspected of being fish-eaters, though some of these reconstructions may be influenced by earlier theories that the group was ancestral to cetaceans.
Phylogeny and evolutionary relationships
Mesonychians were long considered to be creodonts, but have now been removed from that order and placed in three families (Mesonychidae, Hapalodectidae, and Triisodontidae), either within their own order, Mesonychia, or within the order Condylarthra as part of the cohort or superorder Laurasiatheria. Nearly all mesonychians are, on average, larger than most of the Paleocene and Eocene creodonts and miacoid carnivorans.
The order is sometimes referred to by its older name Acreodi.
A recent study found mesonychians to be basal euungulates most closely related to the "arctocyonids" Mimotricentes, Deuterogonodon and Chriacus. "Triisodontidae" may be paraphyletic.
Relationship with whales
Mesonychians possess unusual triangular molar teeth that are similar to those of Cetacea (whales and dolphins), especially those of the archaeocetes, as well as having similar skull anatomies and other morphologic traits. For this reason, scientists had long believed that mesonychians were the direct ancestor of Cetacea, but the discovery of well-preserved hind limbs of archaic cetaceans, as well as more recent phylogenetic analyses now indicate cetaceans are more closely related to hippopotamids and other artiodactyls than they are to mesonychians, and this result is consistent with many molecular studies. The similarity in dentition and skull may be the result of primitive ungulate structures in related groups independently evolving to meet similar needs as predators; some researchers have suggested that the absence of a first toe and a reduced metatarsal are basal features (synapomorphies) indicating that mesonychians, perissodactyls, and artiodactyls are sister groups.
Most paleontologists now doubt that whales are descended from mesonychians, and instead suggest mesonychians are descended from basal ungulates, and that cetaceans are descended from advanced ungulates (Artiodactyla), either deriving from, or sharing a common ancestor with, anthracotheres (the semiaquatic ancestors of hippos). However, the close grouping of whales with hippopotami in cladistic analyses only surfaces following the deletion of Andrewsarchus, which has often been included within the mesonychians. One possible conclusion is that Andrewsarchus has been incorrectly classified. The current uncertainty may, in part, reflect the fragmentary nature of the remains of some crucial fossil taxa, such as Andrewsarchus.
| Biology and health sciences | Mammals: General | Animals |
2007779 | https://en.wikipedia.org/wiki/Water%20bottle | Water bottle | A water bottle is a container that is used to hold liquids, mainly water, for the purpose of transporting a drink while travelling or while otherwise away from a supply of potable water.
Water bottles are usually made of plastic, glass, metal, or some combination of those substances. In the past, water bottles were sometimes made of wood, bark, or animal skins such as leather, hide and sheepskin.
Water bottles can be either disposable or reusable. Disposable water bottles are often sold filled with potable water, while reusable bottles are often sold empty. Reusable water bottles help cut down on consumer plastic waste and carbon emissions.. A reusable water bottle designed for outdoor activities is also called a canteen.
Types
Single-use plastic
Sales of single-use, pre-filled plastic water bottles have increased almost every year for more than a decade. In 2011, greater than US$11 billion was spent on bottled water products in the United States alone. The International Bottled Water Association (IBWA) states that people are increasingly relying on water bottles for convenience and portability.
In some countries with low-quality tap water, citizens also use bottled water (including in family-size containers kept in the home) for health reasons. For example, as of 2010, Mexico had an average 8 percent increase per year in bottled water purchases and consumed approximately 13 percent of the world's total of bottled water. Mexican citizens drink more bottled water than people of any other country, at an average of 61.8 gallons per person each year – more than twice the rate of US per capita consumption. The increase in the use of single-use personal plastic water bottles has contributed markedly to the country's litter problem, though the increase in the popularity of bottled water has come with a decrease in the growth rate of consumption of soft drinks (which pose health risks in excessive quantities, as well as the same littering problem).
Reusable plastic
Multi-use water bottles can be made from high-density polyethylene (HDPE), low-density polyethylene (LDPE), copolyester, or polypropylene. They all offer the advantage of being durable, lightweight, dishwasher-safe, and BPA-free. The main difference between each type of water bottle is the flexibility of the material. Copolyester and polypropylene offer the greatest rigidity; HDPE retains some flexibility; LDPE (most commonly associated with collapsible, squeeze bottles) is highly flexible.
Metal
Metal water bottles are growing in popularity. Made primarily from stainless steel or aluminium (aluminum), they are durable and retain less odor and taste from previous contents than most plastic bottles. But these can sometimes impart a metallic taste. Metal bottles thus often contain a resin or epoxy liner to protect contents from taste and odor transfer or corrosion. Although most liners are now BPA-free, older and less expensive models can contain BPA. Glass liners may also be used .
It is not recommended to fill aluminium bottles with acidic liquids (e.g. orange juice), as this could cause aluminium to leach into the contents of the bottle. Depending on the type of source material and manufacturing process behind a stainless steel bottle, trace amounts of minerals can leach into contents from this type of bottle as well. Stainless steel bottles that do not contain a liner have been known to transfer a rusty taste and odor to contents. Bottles made with food-grade stainless steel (grade 304, also known as 18/8) do not transfer taste or odor.
Metal (especially steel) water bottles can be heavier than their plastic counterparts. Single-walled metal bottles readily transfer temperature of contents to external surfaces, which makes them unsuitable for use with unusually hot or cold liquids. Double-walled metal bottles are insulated to keep cold liquids cold and hot liquids hot, without the external surface being too hot or too cold. Because double-walled bottles have more metal in them, they are more expensive. They are typically vacuum-insulated, but some may have a solid or gel insulation between the metal walls.
Glass
Glass flasks have been used since ancient times, though were not common until the Early Modern period when consistent, bulk manufacturing of glass products became easier. Because they are completely recyclable, BPA-free, and do not retain and transfer taste or odor, glass water bottles are becoming a popular choice for many consumers who are concerned about their health.
Glass bottles are heavier than plastic, stainless steel, or aluminium, and are easier to damage or completely break. Like metal, they also have a high level of temperature transfer, so they are not ideal for very hot or cold liquids. Some types of vacuum-insulated flasks use an inner layer of glass (which is easy to clean), and an outer layer of metal or plastic which helps shield the glass from breakage. Such bottles may still break if dropped, and thus some brands are triple-layer, with the glass inside two layers of plastic; this is a common configuration for large flasks intended for coffee or other liquids that need to be insulated.
Filtering
This type of bottle is often BPA-Free and more commonly uses carbon (activated charcoal) filtration. UV light can also be used to purify water. UV filtration bottles are popular and convenient for those who are travelling to areas where water quality may be harmful, or where bottled water is not readily available. UV is effective against all water-borne pathogens.
Carbon filtration bottles will eliminate some organic chemicals and improve the taste and odor of water. Carbon filtration will not eliminate pathogens, metals or nitrates from water.
Wirelessly connected
Connected devices collect data related to a person's water intake. The data is transmitted to a smartphone, which enables tracking of an individual's water intake and alerts the user when they are not properly hydrated. These devices are a result of technology advancements that fall in the broader category of the Internet of Things. Devices that monitor and collect data related to one's personal health are also part of the quantified self movement. While several concepts have been introduced, none are currently available commercially.
Hydration reservoirs
Hydration reservoirs, also known as hydration bladders, are large-volume, flexible bags typically carried in a backpack system. Users access water via a sipping tube. This system allows the user to remain engaged in activity without having to stop and unscrew a water bottle. Such reservoirs also permit the carrying of a larger water supply (thus a longer hike), as they have both more capacity and better integration into the carrying equipment than an external water bottle or canteen attached to the pack or belt.
Popularity
Due to growing concern over the environmental impact and cost of disposable plastic water bottles, more people are choosing to fill multi-use water bottles. However, the popularity and availability of disposable plastic water bottles continues to rise. In 2007, Americans consumed 50 billion single-serve bottles of water. Since 2001, the sale of single-serve bottled water has fluctuated by 70 percent, and this trend is continuing.
In 2016, a trend among Americans called "water bottle flipping" attracted media attention.
Health
Chemicals used for making some types of bottles have been shown to be detrimental to the health of humans. Inhalation of chemicals used in the manufacture of plastics is a hazard for the factory workers who handle the material. In many developing countries, plastic waste is burned rather than recycled or deposited in landfills. Rural residents of developing countries who burn plastic as a disposal method are not protected from the chemical inhalation hazards associated with this practice. It is important to dispose of water that has been stored in PET bottles that have been exposed to high temperatures for a prolonged period of time or are beyond the expiration date because harmful chemicals may leach from the plastic.
In 2008, researchers from Arizona State University found that storing plastic bottles in temperatures at or above 60 °C can cause antimony to enter the water contained in the bottles. Therefore, frequently drinking from bottles stored in places such as cars during the summer months may have negative health effects.
Bottle manufacturing relies on petroleum and natural resources. Some creating processes release toxic chemicals into the air and water supply that can immensely affect nervous systems, blood cells, kidneys, immune systems, and can cause cancer and birth defects. Most disposable water bottles are made from petroleum found polyethylene terephthalate (PET). While PET is considered less toxic than many other types of plastic, the Berkeley Ecology Center found that manufacturing PET makes toxic emissions in the form of nickel, ethylbenzene, ethylene oxide and benzene at levels 100 times higher than those created to make the same amount of glass.
Environment
Water bottles made of glass, aluminium and steel are the most readily recyclable. HDPE and LDPE bottles can be recycled as well.
Because the manufacturing and transportation of disposable water bottles requires petroleum, a non-renewable resource, the single-serve bottled water industry has come under pressure from concerned consumers. The Pacific Institute calculates that it required about 17 million barrels of oil to make the disposable plastic bottles for single-serve water that Americans consumed in 2006. To sustain the consumptive use of products relying on plastic components and level of manufactured demand for plastic water bottles, the result is shortages of fossil fuels. Furthermore, it means not only a shortage of the raw materials to make plastics, but also a shortage of the energy required to fuel their production.
Single-serve bottled water industry has responded to consumer concern about the environmental impact of disposable water bottles by significantly reducing the amount of plastic used in bottles. The reduced plastic content also results in a lower weight product that uses less energy to transport. Other bottle manufacturing companies are experimenting with alternative materials such as corn starch to make new bottles that are more readily biodegradable.
The lowest impact water bottles are those made of glass or metal. They are not made from petroleum and are easily recyclable. By choosing to continuously fill any multi-use water bottle, the consumer keeps disposable bottles out of the waste stream and minimizes environmental impact.
| Technology | Containers | null |
2008215 | https://en.wikipedia.org/wiki/Quark%20model | Quark model | In particle physics, the quark model is a classification scheme for hadrons in terms of their valence quarks—the quarks and antiquarks that give rise to the quantum numbers of the hadrons. The quark model underlies "flavor SU(3)", or the Eightfold Way, the successful classification scheme organizing the large number of lighter hadrons that were being discovered starting in the 1950s and continuing through the 1960s. It received experimental verification beginning in the late 1960s and is a valid and effective classification of them to date. The model was independently proposed by physicists Murray Gell-Mann, who dubbed them "quarks" in a concise paper, and George Zweig, who suggested "aces" in a longer manuscript. André Petermann also touched upon the central ideas from 1963 to 1965, without as much quantitative substantiation. Today, the model has essentially been absorbed as a component of the established quantum field theory of strong and electroweak particle interactions, dubbed the Standard Model.
Hadrons are not really "elementary", and can be regarded as bound states of their "valence quarks" and antiquarks, which give rise to the quantum numbers of the hadrons. These quantum numbers are labels identifying the hadrons, and are of two kinds. One set comes from the Poincaré symmetry—JPC, where J, P and C stand for the total angular momentum, P-symmetry, and C-symmetry, respectively.
The other set is the flavor quantum numbers such as the isospin, strangeness, charm, and so on. The strong interactions binding the quarks together are insensitive to these quantum numbers, so variation of them leads to systematic mass and coupling relationships among the hadrons in the same flavor multiplet.
All quarks are assigned a baryon number of . Up, charm and top quarks have an electric charge of +, while the down, strange, and bottom quarks have an electric charge of −. Antiquarks have the opposite quantum numbers. Quarks are spin- particles, and thus fermions. Each quark or antiquark obeys the Gell-Mann–Nishijima formula individually, so any additive assembly of them will as well.
Mesons are made of a valence quark–antiquark pair (thus have a baryon number of 0), while baryons are made of three quarks (thus have a baryon number of 1). This article discusses the quark model for the up, down, and strange flavors of quark (which form an approximate flavor SU(3) symmetry). There are generalizations to larger number of flavors.
History
Developing classification schemes for hadrons became a timely question after new experimental techniques uncovered so many of them that it became clear that they could not all be elementary. These discoveries led Wolfgang Pauli to exclaim "Had I foreseen that, I would have gone into botany." and Enrico Fermi to advise his student Leon Lederman: "Young man, if I could remember the names of these particles, I would have been a botanist." These new schemes earned Nobel prizes for experimental particle physicists, including Luis Alvarez, who was at the forefront of many of these developments. Constructing hadrons as bound states of fewer constituents would thus organize the "zoo" at hand. Several early proposals, such as the ones by Enrico Fermi and Chen-Ning Yang (1949), and the Sakata model (1956), ended up satisfactorily covering the mesons, but failed with baryons, and so were unable to explain all the data.
The Gell-Mann–Nishijima formula, developed by Murray Gell-Mann and Kazuhiko Nishijima, led to the Eightfold Way classification, invented by Gell-Mann, with important independent contributions from Yuval Ne'eman, in 1961. The hadrons were organized into SU(3) representation multiplets, octets and decuplets, of roughly the same mass, due to the strong interactions; and smaller mass differences linked to the flavor quantum numbers, invisible to the strong interactions. The Gell-Mann–Okubo mass formula systematized the quantification of these small mass differences among members of a hadronic multiplet, controlled by the explicit symmetry breaking of SU(3).
The spin- baryon, a member of the ground-state decuplet, was a crucial prediction of that classification. After it was discovered in an experiment at Brookhaven National Laboratory, Gell-Mann received a Nobel Prize in Physics for his work on the Eightfold Way, in 1969.
Finally, in 1964, Gell-Mann and George Zweig, discerned independently what the Eightfold Way picture encodes: They posited three elementary fermionic constituents—the "up", "down", and "strange" quarks—which are unobserved, and possibly unobservable in a free form. Simple pairwise or triplet combinations of these three constituents and their antiparticles underlie and elegantly encode the Eightfold Way classification, in an economical, tight structure, resulting in further simplicity. Hadronic mass differences were now linked to the different masses of the constituent quarks.
It would take about a decade for the unexpected nature—and physical reality—of these quarks to be appreciated more fully (See Quarks). Counter-intuitively, they cannot ever be observed in isolation (color confinement), but instead always combine with other quarks to form full hadrons, which then furnish ample indirect information on the trapped quarks themselves. Conversely, the quarks serve in the definition of quantum chromodynamics, the fundamental theory fully describing the strong interactions; and the Eightfold Way is now understood to be a consequence of the flavor symmetry structure of the lightest three of them.
Mesons
The Eightfold Way classification is named after the following fact: If we take three flavors of quarks, then the quarks lie in the fundamental representation, 3 (called the triplet) of flavor SU(3). The antiquarks lie in the complex conjugate representation . The nine states (nonet) made out of a pair can be decomposed into the trivial representation, 1 (called the singlet), and the adjoint representation, 8 (called the octet). The notation for this decomposition is
Figure 1 shows the application of this decomposition to the mesons. If the flavor symmetry were exact (as in the limit that only the strong interactions operate, but the electroweak interactions are notionally switched off), then all nine mesons would have the same mass. However, the physical content of the full theory includes consideration of the symmetry breaking induced by the quark mass differences, and considerations of mixing between various multiplets (such as the octet and the singlet).
N.B. Nevertheless, the mass splitting between the and the is larger than the quark model can accommodate, and this "– puzzle" has its origin in topological peculiarities of the strong interaction vacuum, such as instanton configurations.
Mesons are hadrons with zero baryon number. If the quark–antiquark pair are in an orbital angular momentum state, and have spin , then
≤ J ≤ L + S, where S = 0 or 1,
P = (−1)L+1, where the 1 in the exponent arises from the intrinsic parity of the quark–antiquark pair.
C = (−1)L+S for mesons which have no flavor. Flavored mesons have indefinite value of C.
For isospin and 0 states, one can define a new multiplicative quantum number called the G-parity such that .
If P = (−1)J, then it follows that S = 1, thus PC = 1. States with these quantum numbers are called natural parity states; while all other quantum numbers are thus called exotic (for example, the state ).
Baryons
Since quarks are fermions, the spin–statistics theorem implies that the wavefunction of a baryon must be antisymmetric under the exchange of any two quarks. This antisymmetric wavefunction is obtained by making it fully antisymmetric in color, discussed below, and symmetric in flavor, spin and space put together. With three flavors, the decomposition in flavor is
The decuplet is symmetric in flavor, the singlet antisymmetric and the two octets have mixed symmetry. The space and spin parts of the states are thereby fixed once the orbital angular momentum is given.
It is sometimes useful to think of the basis states of quarks as the six states of three flavors and two spins per flavor. This approximate symmetry is called spin-flavor SU(6). In terms of this, the decomposition is
The 56 states with symmetric combination of spin and flavour decompose under flavor SU(3) into
where the superscript denotes the spin, S, of the baryon. Since these states are symmetric in spin and flavor, they should also be symmetric in space—a condition that is easily satisfied by making the orbital angular momentum . These are the ground-state baryons.
The octet baryons are the two nucleons (, ), the three Sigmas (, , ), the two Xis (, ), and the Lambda (). The decuplet baryons are the four Deltas (, , , ), three Sigmas (, , ), two Xis (, ), and the Omega ().
For example, the constituent quark model wavefunction for the proton is
Mixing of baryons, mass splittings within and between multiplets, and magnetic moments are some of the other quantities that the model predicts successfully.
The group theory approach described above assumes that the quarks are eight components of a single particle, so the anti-symmetrization applies to all the quarks. A simpler approach is to consider the eight flavored quarks as eight separate, distinguishable, non-identical particles. Then the anti-symmetrization applies only to two identical quarks (like uu, for instance).
Then, the proton wavefunction can be written in a simpler form:
and the
If quark–quark interactions are limited to two-body interactions, then all the successful quark model predictions, including sum rules for baryon masses and magnetic moments, can be derived.
Discovery of color
Color quantum numbers are the characteristic charges of the strong force, and are completely uninvolved in electroweak interactions. They were discovered as a consequence of the quark model classification, when it was appreciated that the spin baryon, the , required three up quarks with parallel spins and vanishing orbital angular momentum. Therefore, it could not have an antisymmetric wavefunction, (required by the Pauli exclusion principle). Oscar Greenberg noted this problem in 1964, suggesting that quarks should be para-fermions.
Instead, six months later, Moo-Young Han and Yoichiro Nambu suggested the existence of a hidden degree of freedom, they labeled as the group SU(3)' (but later called 'color). This led to three triplets of quarks whose wavefunction was anti-symmetric in the color degree of freedom.
Flavor and color were intertwined in that model: they did not commute.
The modern concept of color completely commuting with all other charges and providing the strong force charge was articulated in 1973, by William Bardeen, Harald Fritzsch, and Murray Gell-Mann.
States outside the quark model
While the quark model is derivable from the theory of quantum chromodynamics, the structure of hadrons is more complicated than this model allows. The full quantum mechanical wavefunction of any hadron must include virtual quark pairs as well as virtual gluons, and allows for a variety of mixings. There may be hadrons which lie outside the quark model. Among these are the glueballs (which contain only valence gluons), hybrids (which contain valence quarks as well as gluons) and exotic hadrons (such as tetraquarks or pentaquarks).
| Physical sciences | Particle physics: General | Physics |
2008686 | https://en.wikipedia.org/wiki/Poikilotherm | Poikilotherm | A poikilotherm () is an animal (Greek poikilos – 'various', 'spotted', and therme – 'heat') whose internal temperature varies considerably. Poikilotherms have to survive and adapt to environmental stress. One of the most important stressors is outer environment temperature change, which can lead to alterations in membrane lipid order and can cause protein unfolding and denaturation at elevated temperatures. Poikilotherm is the opposite of homeotherm – an animal which maintains thermal homeostasis. In principle, the term could be applied to any organism, but it is generally only applied to vertebrate animals. Usually the fluctuations are a consequence of variation in the ambient environmental temperature. Many terrestrial ectotherms are poikilothermic. However some ectotherms seek constant-temperature environments to the point that they are able to maintain a constant internal temperature, and are considered actual or practical homeotherms. It is this distinction that often makes the term poikilotherm more useful than the vernacular "cold-blooded", which is sometimes used to refer to ectotherms more generally.
Poikilothermic animals include types of vertebrate animals, specifically some fish, amphibians, and reptiles, as well as many invertebrate animals. The naked mole-rat and sloths are some of the rare mammals which are poikilothermic.
Etymology
The term derives from Greek poikilos (), meaning "varied," ultimately from a root meaning "dappled" or "painted," and thermos (), meaning "heat".
Physiology
Poikilotherm animals must be able to function over a wider range of temperatures than homeotherms. The speed of most chemical reactions vary with temperature, and in order to function poikilotherms may have four to ten enzyme systems that operate at different temperatures for an important chemical reaction. As a result, poikilotherms often have larger, more complex genomes than homeotherms in the same ecological niche. Frogs are a notable example of this effect, though their complex development is also an important factor in their large genome.
Because their metabolism is variable and generally below that of homeothermic animals, sustained high-energy activities like powered flight in large animals or maintaining a large brain is generally beyond poikilotherm animals. The metabolism of poikilotherms favors strategies such as sit-and-wait hunting over chasing prey for larger animals with high movement cost. As they do not use their metabolisms to heat or cool themselves, total energy requirement over time is low. For the same body weight, poikilotherms need only 5 to 10% of the energy of homeotherms.
Adaptations in poikilotherms
Some adaptations are behavioral. Lizards and snakes bask in the sun in the early morning and late evening, and seek shelter around noon.
The eggs of the yellow-faced bumblebee are unable to regulate heat. A behavioral adaptation to combat this is incubation, where to maintain the internal temperatures of eggs, the queen and her workers will incubate the brood almost constantly, by warming their abdomens and touching them to the eggs. The bumblebee generates heat by shivering flight muscles even though it is not flying.
Termite mounds are usually oriented in a north–south direction so that they absorb as much heat as possible around dawn and dusk and minimise heat absorption around noon.
Tuna are able to warm their entire bodies through a heat exchange mechanism called the rete mirabile, which helps keep heat inside the body, and minimises the loss of heat through the gills. They also have their swimming muscles near the center of their bodies instead of near the surface, which minimises heat loss.
Gigantothermy means growing to large size in order to reduce heat loss, such as in sea turtles and ice-age megafauna. Body volume increases proportionally faster than does body surface, with increasing size; and less body surface area per unit body volume tends to minimise heat loss.
Camels, although they are homeotherms, thermoregulate using a method termed "temperature cycling" to conserve energy. In hot deserts, they allow their body temperature to rise during the day and fall during the night, adjusting their body temperature to cycle over approximately 6 °C.
Ecology
It is comparatively easy for a poikilotherm to accumulate enough energy to reproduce. Poikilotherms at the same trophic level often have much shorter generations than homeotherms: weeks rather than years. Such applies even to animals with similar ecological roles such as cats and snakes.
This difference in energy requirement also means that a given food source can support a greater density of poikilothermic animals than homeothermic animals. This is reflected in the predator-prey ratio which is usually higher in poikilothermic fauna compared to homeothermic ones. However, when homeotherms and poikilotherms have similar niches, and compete, the homeotherm can often drive poikilothermic competitors to extinction, because homeotherms can gather food for a greater fraction of each day and in more effective, specialized ways (e.g. chimpanzees actively seeking out and collecting army ants with sticks versus the typical poikilotherm sit-and-wait strategy).
In medicine
In medicine, loss of normal thermoregulation is referred to as poikilothermia. This can be seen in compartment syndrome and with use of sedative-hypnotics like barbiturates, ethanol, and chloral hydrate. REM sleep is considered a poikilothermic state in humans. Poikilothermia is one of the signs of acute limb ischemia.
| Biology and health sciences | Basics | Biology |
1370818 | https://en.wikipedia.org/wiki/Blackbuck | Blackbuck | The blackbuck (Antilope cervicapra), also known as the Indian antelope, is a medium-sized antelope native to India and Nepal. It inhabits grassy plains and lightly forested areas with perennial water sources.
It stands up to high at the shoulder. Males weigh , with an average of . Females are lighter, weighing or on average. Males have long corkscrew horns, and females occasionally develop horns, as well. The white fur on the chin and around the eyes is in sharp contrast with the black stripes on the face. Both sexes' coats feature a two-tone colouration; in males, the majority of the body is dark brown to black, with white circles around the eyes, white ears and tail, and the belly, lower jaw, and inner legs also white. Females and juveniles are yellowish-fawn to tan and display the same white areas, only with more of a beige tone than the males. Females also feature a more pronounced horizontal white side-stripe, starting around the shoulder and ending at the rump. The blackbuck is the sole living member of the genus Antilope and was described by Carl Linnaeus in 1758. Two subspecies are recognized.
The blackbuck is active mainly during the day. It forms three types of small groups: female, male, and young bachelor herds. Males often adopt lekking as a strategy to garner females for mating. While other males are not allowed into these territories, females often visit these places to forage. The male can thus attempt mating with her. The blackbuck is an herbivore and grazes on low grasses, occasionally browsing as well. Females become sexually mature at the age of eight months, but mate no earlier than two years of age. Males mature later, at 1.5 years. Mating takes place throughout the year. Gestation is typically six months long, after which a single calf is born. The lifespan is typically 10 to 15 years.
The antelope is native to and occurs mainly in India, while it is locally extinct in Pakistan and Bangladesh. Formerly widespread, small and scattered herds are largely confined to protected areas today. During the 20th century, blackbuck numbers declined sharply due to excessive hunting, deforestation, and habitat destruction. The blackbuck has been introduced in Argentina, Australia and the United States, primarily on hunting ranches. In Argentina, the population is surviving well. In India, hunting of blackbuck is prohibited under Schedule I of the Wildlife Protection Act of 1972. The blackbuck has significance in Hinduism; Indian and Nepali villagers do not harm the antelope.
Etymology
The scientific name of the blackbuck Antilope cervicapra stems from the Latin word antalopus ("horned animal"). The specific name cervicapra is composed of the Latin words cervus ("deer") and capra ("she-goat"). The vernacular name "blackbuck" is a reference to the dark brown to black colour of the dorsal part of the coat of the males. The earliest recorded use of this name dates back to 1850.
Taxonomy and evolution
The blackbuck is the sole living member of the genus Antilope and is classified in the family Bovidae. The species was described and given its binomial name by Swedish zoologist Carl Linnaeus in the 10th edition of Systema Naturae in 1758. Antilope also includes fossil species, such as Antilope subtorta and Antilope intermedia.
Antilope, Eudorcas, Gazella, and Nanger form a clade within their tribe Antilopini. A 1995 study of the detailed karyotype of Antilope suggested that within this clade, Antilope is closest to the Gazella group. A 1999 phylogenetic analysis confirmed that Antilope is the closest sister taxon to Gazella, although an earlier phylogeny, proposed in 1976, placed Antilope as sister to Nanger. In a more recent revision of the phylogeny of the Antilopini on the basis of sequences from multiple nuclear and mitochondrial loci in 2013, Eva Verena Bärmann (of the University of Cambridge) and colleagues re-examined the phylogenetic relationships and found Antilope and Gazella to be sister genera distinct from the sister genera Nanger and Eudorcas.
Two subspecies are recognised, although they might be independent species:
A. c. cervicapra (Linnaeus, 1758), known as the southeastern blackbuck, occurs in southern, eastern, and central India. The white eye ring of the male is narrow above the eye and the neck is all black in the male and the white on the underside is largely restricted to the belly in both males and females. The black leg stripe is well defined and reaches all along the leg.
A. c. rajputanae Zukowsky, 1927, known as the northwestern blackbuck, occurs in northwestern India. Males have a grey sheen to the dark parts during the breeding season. The white on the underside extends up to half way on the sides of the body and the lower base of the neck of males is white. The white eye ring is broad all around the eye with the leg-stripe going only down to the shanks.
Genetics
The blackbuck shows variation in its diploid chromosome number. Males have 31–33, while females have 30–32. Males have an XY1Y2 sex chromosome. Unusually large sex chromosomes had earlier been described only in a few species, all of which belonged to Rodentia. However, in 1968, a study found that two artiodactyls, the blackbuck and the sitatunga, too, showed this abnormality. Generally, the X chromosome constitutes 5% of the haploid chromosomal complement, but the X chromosome of the blackbuck this percentage is 14.96. Portions of both peculiarly large chromosomes show delayed replication.
A 1997 study found lower variation in blood protein polymorphism in Antilope in comparison with Antidorcas, Eudorcas, and Gazella. This was attributed to a history of rapid evolution of an autapomorphic phenotype of Antilope. This might have been aided by a particularly strong selection of a few dominant males due to their lekking behaviour.
Characteristics
The blackbuck has white fur on the chin and around the eyes, which is in sharp contrast with the black stripes on the face. The coats of males show two-tone colouration; while the upper parts and outsides of the legs are dark brown to black, the underparts and the insides of the legs are all white. Darkness typically increases as the male ages; females and juveniles are yellowish fawn to tan. In Texas, blackbuck moult in spring, following which the males look notably lighter, though darkness persists on the face and the legs. On the contrary, males grow darker as the breeding season approaches. Both melanism and albinism have been observed in wild blackbuck. Albino blackbuck are often zoo attractions as in the Indira Gandhi Zoological Park.
The blackbuck is a moderately sized antelope. It stands up to high at the shoulder; the head-to-body length is nearly . In the population introduced to Texas, males weigh , an average of . Females are lighter, weighing or on average. Sexual dimorphism is prominent, as males are heavier and darker than the females. The long, ringed horns, that resemble corkscrews, are generally present only on males, though females may also develop horns. They measure , though the maximum horn length recorded in Texas has not exceeded . The horns diverge forming a "V"-like shape. In India, horns are longer and more divergent in specimens from the northern and western parts of the country.
Blackbuck bear a close resemblance to gazelles, and are distinguished mainly by the fact that while gazelles are brown in the dorsal parts, blackbuck develop a dark brown or black colour in these parts.
Distribution and habitat
The blackbuck is native to the Indian subcontinent and inhabits grassy plains and thinly forested areas where perennial water sources are available for its daily need to drink. Herds travel long distances to obtain water.
The British naturalist William Thomas Blanford described the range of the blackbuck in his 1891 The Fauna of British India, including Ceylon and Burma as:
Today, small, scattered herds are largely confined to protected areas.
In Pakistan, the blackbuck occasionally occurred along the border with India until 2001. In southern Nepal, the last surviving blackbuck population in Blackbuck Conservation Area was estimated to comprise 184 individuals in 2008.
A few blackbucks are present in the Indian Institute of Technology Madras campus.
The blackbuck is considered locally extinct in Pakistan and Bangladesh.
Introduced populations
The blackbuck was also introduced into Argentina, numbering about 8,600 individuals as of the early 2000s.
In the early 1900s, blackbuck were introduced to Western Australia.
In either the late 1980s or the early 1990s, they were also introduced to Cape York in Far North Queensland, although the population was subsequently eradicated.
In 2013, an antelope that appeared to be a blackbuck was sighted at Kakadu National Park in the Northern Territory.
In 2015, a blackbuck was sighted near Warrnambool, Victoria, which was later captured and sent to Mansfield Zoo.
The blackbuck is a declared pest in Queensland and Western Australia. In Victoria, blackbuck and American bison are considered both "regulated pest animals" and livestock.
The antelope was introduced in Texas in the Edwards Plateau in 1932. By 1988, the population had increased and the antelope was the most populous exotic animal in Texas after the chital.
Ecology and behaviour
The blackbuck is a diurnal antelope, though is less active at noon when summer temperatures rise. It can run at a speed of .
Group size fluctuates and seems to depend on the availability of forage and the nature of the habitat. Large herds have an edge over smaller ones in that danger can be detected faster, though individual vigilance is lower in the former. Large herds spend more time feeding than small herds. A disadvantage for large herds, however, is that traveling requires more resources. Herd size reduces in summer.
Males often adopt lekking as a strategy on the part of males to garner females for mating. Territories are established by males on the basis of the local distribution of female groups, which in turn is determined by the habitat, so as to ensure greater access to females. The males actively defend resources in their territories, nearly in size; territories are marked with scent using preorbital gland and interdigital gland secretions, faeces and urine. While other males are not allowed into these territories, females are allowed to visit these places to forage. The male can attempt mating with visiting females. Lekking is a demanding strategy, as the males often have to bear injuries – thus it is a tactic typically adopted by strong, dominant males. Males may either defend their mates or try to forcibly copulate with them. Weaker males, who may not be dominant, might choose the second method.
The blackbuck is severely affected by natural calamities such as floods and droughts, from which it can take as long as five years to recover. The wolf is a major predator. Old rutting bulls might be especially vulnerable prey. The golden jackal hunts juveniles. Village dogs are reported to kill fawns, but are unlikely to successfully hunt and kill adults.
Blackbucks in Great Indian Bustard Sanctuary show flexible habitat use as the resources and risks change seasonally in the landscape. They use small patches in the area of about . Human activities strongly influenced the movement of herds, but the presence of small refuges allowed them to persist in the landscape.
Diet
The blackbuck is a herbivore and grazes on low grasses, occasionally browsing as well. It prefers sedges, fall witchgrass, mesquite, and live oak and was occasionally observed browsing on acacia trees in the Cholistan Desert. Oats and berseem were found to be palatable and nutritious to captive populations. In Velavadar Black Buck Sanctuary, Dichanthium annulatum comprised 35% of the diet. Digestion of nutrients, especially crude proteins, was poor in summer, but more efficient in the rainy and winter seasons. Crude protein intake in summer was very low, even below the recommended value. Blackbuck consumed less food in summer than in winter, and often foraged on the fruits of Prosopis juliflora. Prosopis becomes a significant food item if grasses are scarce. Water is a daily requirement of the blackbuck.
Reproduction
Females become sexually mature at the age of eight months, but mate no earlier than two years. Males mature at the age of one-and-a-half years. Mating takes place throughout the year; peaks occur during spring and fall in Texas. Two peaks have been observed in India: from August to October and from March to April. Rutting males aggressively establish and defend their territories from other males, giving out loud grunts and engaging in serious head-to-head fights, pushing each other using horns. Aggressive display consists of thrusting the neck forward and raising it, folding the ears and raising the tail. The dominant male pursues the female with his nose pointing upward, smells her urine and shows a flehmen response. The female shows her receptivity by waving her tail and thumping the hindlegs on the ground. This is followed by several mounting attempts, and copulation. The whole process may last as long as six hours. The female will remain still for some time after copulation, following which she may start grazing. The male may then move on to mate with another female.
Gestation typically lasts six months, after which a single calf is born. Newborns are a light yellow; infant males may have a black patch on the head and the neck. Young are precocial, they can stand on their own soon after birth. Females can mate again after a month of parturition. Juveniles remain active and playful throughout the day. Juvenile males turn black gradually, darkening notably after the third year. The lifespan is typically 10 to 15 years.
Threats
During the 20th century, blackbuck numbers declined sharply due to excessive hunting, deforestation and habitat degradation. Some blackbucks are killed illegally especially where the species is sympatric with nilgai.
Until India's independence in 1947, blackbuck and chinkara were hunted in many princely states with specially trained captive Asiatic cheetahs. By the 1970s, blackbuck was locally extinct in several areas.
Conservation
The blackbuck is listed under Appendix III of CITES.
In India, hunting of blackbuck is prohibited under Schedule I of the Wildlife Protection Act of 1972. It inhabits several protected areas of India, including
in Gujarat: Velavadar National Park, Gir Forest National Park;
in Bihar: Kaimur Wildlife Sanctuary;
in Maharashtra: Great Indian Bustard Sanctuary;
in Madhya Pradesh: Kanha National Park
in Rajasthan: Tal Chhapar Sanctuary, National Chambal Sanctuary, Ranthambhore National Park
in Karnataka: Ranibennur Blackbuck Sanctuary;
in Tamil Nadu: Point Calimere Wildlife and Bird Sanctuary, Vallanadu Wildlife Sanctuary, Guindy National Park.
in Punjab: Abohar Wildlife Sanctuary
A captive population is maintained in Pakistan's Lal Suhanra National Park.
In culture
The blackbuck has associations with the Indian culture. The antelope might have been a source of food in the Indus Valley civilisation (3300–1700 BCE); bone remains have been discovered in sites such as Dholavira and Mehrgarh. The blackbuck is routinely depicted in miniature paintings of the Mughal era of 16th to 19th centuries depicting royal hunts often using cheetahs. Villagers in India and Nepal generally do not harm the blackbuck. Tribes such as the Bishnois revere and care for most animals including the blackbuck.
The blackbuck is mentioned in Sanskrit texts as the Kṛṣṇamṛga. According to Hindu mythology, it draws the chariot of Lord Krishna. The blackbuck is considered to be the vehicle of the wind god Vayu, the divine drink Soma and the moon god Chandra. In Tamil Nadu, the blackbuck is considered to be the vehicle of the Hindu goddess Korravai. In Rajasthan, the goddess Karni Mata is believed to protect the blackbuck.
In the Yājñavalkya Smṛti, Sage Yagyavalkya is quoted stating "in what country there is black antelope, in that Dharma must be known", which is interpreted to mean that certain religious practices including sacrifices were not to be performed where blackbuck did not roam.
The hide of the blackbuck is deemed to be sacred in Hinduism. According to the scriptures, it is to be sat upon only by brahmin priests, sadhus and yogis, forest-dwellers and bhikshu mendicants. Blackbuck meat is highly regarded in Texas. In an analysis, blackbuck milk was found to have 6.9% protein, 9.3% fat, and 4.3% lactose.
In some agricultural areas in northern India, the blackbuck are found in large numbers and raid crop fields. However, the damage caused by blackbuck is far lower than that caused by the nilgai.
In 2018, Bollywood actor Salman Khan was sentenced to five years imprisonment for poaching a blackbuck in 1998.
| Biology and health sciences | Bovidae | Animals |
1371050 | https://en.wikipedia.org/wiki/Medicalization | Medicalization | Medicalization is the process by which human conditions and problems come to be defined and treated as medical conditions, and thus become the subject of medical study, diagnosis, prevention, or treatment. Medicalization can be driven by new evidence or hypotheses about conditions; by changing social attitudes or economic considerations; or by the development of new medications or treatments.
Medicalization is studied from a sociologic perspective in terms of the role and power of professionals, patients, and corporations, and also for its implications for ordinary people whose self-identity and life decisions may depend on the prevailing concepts of health and illness. Once a condition is classified as medical, a medical model of disability tends to be used in place of a social model. Medicalization may also be termed pathologization or (pejoratively) "disease mongering". Since medicalization is the social process through which a condition becomes seen as a medical disease in need of treatment, appropriate medicalization may be viewed as a benefit to human society. The identification of a condition as a disease can lead to the treatment of certain symptoms and conditions, which will improve overall quality of life.
History
The concept of medicalization was devised by sociologists to explain how medical knowledge is applied to behaviors which are not self-evidently medical or biological. The term medicalization entered the sociology literature in the 1970s in the works of Irving Zola, Peter Conrad and Thomas Szasz, among others. According to Eric Cassell's book, The Nature of Suffering and the Goals of Medicine (2004), the expansion of medical social control is being justified as a means of explaining deviance. These sociologists viewed medicalization as a form of social control in which medical authority expanded into domains of everyday existence, and they rejected medicalization in the name of liberation. This critique was embodied in works such as Conrad's article "The discovery of hyperkinesis: notes on medicalization of deviance", published in 1973 (hyperkinesis was the term then used to describe what we might now call ADHD).
These sociologists did not believe medicalization to be a new phenomenon, arguing that medical authorities had always been concerned with social behavior and traditionally functioned as agents of social control (Foucault, 1965; Szasz,1970; Rosen). However, these authors took the view that increasingly sophisticated technology had extended the potential reach of medicalization as a form of social control, especially in terms of "psychotechnology" (Chorover,1973).
In the 1975 book Limits to medicine: Medical nemesis (1975), Ivan Illich put forth one of the earliest uses of the term "medicalization". Illich, a philosopher, argued that the medical profession harms people through iatrogenesis, a process in which illness and social problems increase due to medical intervention. Illich saw iatrogenesis occurring on three levels: the clinical, involving serious side effects worse than the original condition; the social, whereby the general public is made docile and reliant on the medical profession to cope with life in their society; and the structural, whereby the idea of aging and dying as medical illnesses effectively "medicalized" human life and left individuals and societies less able to deal with these "natural" processes.
The concept of medicalization dovetailed with some aspects of the 1970s feminist movement. Critics such as Ehrenreich and English (1978) argued that women's bodies were being medicalized by the predominantly male medical profession. Menstruation and pregnancy had come to be seen as medical problems requiring interventions such as hysterectomies.
Marxists such as Vicente Navarro (1980) linked medicalization to an oppressive capitalist society. They argued that medicine disguised the underlying causes of disease, such as social inequality and poverty, and instead presented health as an individual issue. Others examined the power and prestige of the medical profession, including the use of terminology to mystify and of professional rules to exclude or subordinate others.
Tiago Correia (2017) offers an alternative perspective on medicalization. He argues that medicalization needs to be detached from biomedicine to overcome much of the criticism it has faced, and to protect its value in contemporary sociological debates. Building on Gadamer's hermeneutical view of medicine, he focuses on medicine's common traits, regardless of empirical differences in both time and space. Medicalization and social control are viewed as distinct analytical dimensions that in practice may or may not overlap. Correia contends that the idea of "making things medical" needs to include all forms of medical knowledge in a global society, not simply those forms linked to the established (bio)medical professions. Looking at "knowledge", beyond the confines of professional boundaries, may help us understand the multiplicity of ways in which medicalization can exist in different times and societies, and allow contemporary societies to avoid such pitfalls as "demedicalization" (through a turn towards complementary and alternative medicine) on the one hand, or the over-rapid and unregulated adoption of biomedical medicine in non-western societies on the other. The challenge is to determine what medical knowledge is present, and how it is being used to medicalize behaviors and symptoms.
Areas
Sexuality and gender
Many aspects of human sexuality have been medicalized and pathologised by psychiatry, psychology and the pharmaceutical industry. This includes masturbation, homosexuality, erectile dysfunction and female sexual dysfunction. Medicalization has also been used to justify sexualisation of transgender people, intersex people and those diagnosed with HIV/AIDS. The medicalization of sexuality has resulted in increased social control, disease mongering, surveillance, and increased funding in some research areas of sexology and human physiology. The practice of medicalizing sexuality has been widely criticized, with one of the most common criticisms being that the biological reductionism and other tenets of medicalisation, individualism and naturalism, generally fail to take into account sociocultural factors contributing to human sexuality.
The HIV/AIDS pandemic allegedly caused from the 1980s a "profound re-medicalization of sexuality".
The diagnosis of premenstrual dysphoric disorder (PMDD) has caused some controversy when fluoxetine (also known as Prozac) was being repackaged as a PMDD therapy under the trade named Sarafem. The psychologist Peggy Kleinplatz has criticized the diagnosis as the medicalization of normal human behavior. Other medicalized aspects of women's health include infertility, breastfeeding, the childbirth process, and postpartum depression.
Although it has received less attention, it is claimed that masculinity has also faced medicalization, being deemed damaging to health and requiring regulation or enhancement through drugs, technologies or therapy. Specifically, erectile dysfunction was once considered a natural part of the aging process in men, but has since been medicalized as a problem, late-onset hypogonadism.
According to Mike Fitzpatrick, resistance to medicalization was a common theme of the gay liberation, anti-psychiatry, and feminist movements of the 1970s, but now there is "virtually no resistance to the advance of government intrusion in lifestyle if it is deemed to be justified in terms of public health." Moreover, the pressure for medicalization now comes from society itself as well as from the government and medical professionals.
Psychiatry
For many years, marginalized psychiatrists (such as Peter Breggin, Paula Caplan, Thomas Szasz) and outside critics (such as Stuart A. Kirk) have "been accusing psychiatry of engaging in the systematic medicalization of normality". More recently these concerns have come from insiders who have worked for and promoted the American Psychiatric Association (e.g., Robert Spitzer, Allen Frances).
Benjamin Rush, the father of American psychiatry, claimed that Black people had black skin because they were ill with hereditary leprosy. Consequently, he considered vitiligo as a "spontaneous cure".
According to Franco Basaglia and his followers, whose approach pointed out the role of psychiatric institutions in the control and medicalization of deviant behaviors and social problems, psychiatry is used as the provider of scientific support for social control to the existing establishment, and the ensuing standards of deviance and normality brought about repressive views of discrete social groups. As scholars have long argued, governmental and medical institutions code menaces to authority as mental diseases during political disturbances.
According to Nicholas Kittrie, a number of phenomena considered "deviant", such as alcoholism, drug addiction, prostitution, pedophilia, and masturbation ("self-abuse"), were originally considered as moral, then legal, and now medical problems. Innumerable other conditions such as obesity, smoking cigarettes, draft malingering, bachelorhood, divorce, unwanted pregnancy, kleptomania, and grief, have been declared diseases by medical and psychiatric authorities. Due to these perceptions, peculiar deviants were subjected to moral, then legal, and now medical modes of social control. Similarly, Conrad and Schneider concluded their review of the medicalization of deviance by identifying three major paradigms that have reigned over deviance designations in different historical periods: deviance as sin; deviance as crime; and deviance as sickness.
According to Thomas Szasz, "the therapeutic state swallows up everything human on the seemingly rational ground that nothing falls outside the province of health and medicine, just as the theological state had swallowed up everything human on the perfectly rational ground that nothing falls outside the province of God and religion".
Labeling theory
A 2002 editorial in the British Medical Journal warned of inappropriate medicalization leading to disease mongering, where the boundaries of the definition of illnesses are expanded to include personal problems as medical problems or risks of diseases are emphasized to broaden the market for medications. The authors noted:
Healthism
Public health campaigns have been criticized as a form of "healthism", which is moralistic in nature rather than primarily focused on health. Medical doctors Petr Shkrabanek and James McCormick wrote a series of publications on this topic in the late 1980s and early 1990s criticizing the UK's Health of The Nation campaign. These publications exposed abuse of epidemiology and statistics by public health authorities and organizations to support lifestyle interventions and screening programs. Inculcating a fear of ill-health and a strong notion of individual responsibility has been derided as "health fascism" by some scholars as it objectifies the individual without considering emotional or social factors.
Professionals, patients, corporations and society
Several decades on the definition of medicalization is complicated, if for no other reason than because the term is so widely used. Many contemporary critics position pharmaceutical companies in the space once held by doctors as the supposed catalysts of medicalization. Titles such as "The making of a disease" or "Sex, drugs, and marketing" critique the pharmaceutical industry for shunting everyday problems into the domain of professional biomedicine. At the same time, others reject as implausible any suggestion that society rejects drugs or drug companies and highlight that the same drugs that are allegedly used to treat deviances from societal norms also help many people live their lives. Even scholars who critique the societal implications of brand-name drugs generally remain open to these drugs' curative effects – a far cry from earlier calls for a revolution against the biomedical establishment. The emphasis in many quarters has come to be on "overmedicalization" rather than "medicalization" in itself.
Others, however, argue that in practice the process of medicalization tends to strip subjects of their social context, so they come to be understood in terms of the prevailing biomedical ideology, resulting in a disregard for overarching social causes such as unequal distribution of power and resources. A series of publications by Mens Sana Monographs have focused on medicine as a corporate capitalist enterprise.
Scholars argue that in the late 20th century transformation within the health sector in the US altered the relationship between people in the healthcare sector. This has been attributed to the commodification of healthcare and the role of parties other than doctors such as insurance companies, the pharmaceutical industry, and the government, referred to collectively as countervailing powers. The doctor remains an authority figure who prescribes pharmaceuticals to patients. However, in some countries, such as the US, ubiquitous direct-to-consumer advertising encourages patients to ask for particular drugs by name, thereby creating a conversation between consumer and drug company that threatens to cut the doctor out of the loop. Additionally, there is a widespread concern regarding the extent of the pharmaceutical marketing direct to doctors and other healthcare professionals. Examples of this direct marketing are visits by salespeople, funding of journals, training courses or conferences, incentives for prescribing, and the routine provision of "information" written by the pharmaceutical company. The role of patients in this economy has also changed. Once regarded as passive victims of medicalization, patients can now occupy active positions as advocates, consumers, or even agents of change.
In response to theory based on medicalisation being insufficient to explain social processes, some scholars have developed a concept of biomedicalization which argues that technical and scientific interventions are transforming medicine. One aspect is pharmaceuticalization, the influence of the use of pharmaceutical drugs rather than other interventions. Other components are computerization of parts of healthcare such as public health, the creation of a "biopolitical economy" of private research outside of state, the perception of health as a moral obligation.
Medicalization has brought health issues to the fore, so people think more and more about things in terms of health and act to promote health. When it comes to health issues, medicine is not the only provider of answers, but there have always been alternatives and competitors. At the same time as medicalization, "paramedicalization" has strengthened: also many treatments for which there is no medical basis, at least for now, are popular and commercially successful.
| Biology and health sciences | General concepts | Health |
1372353 | https://en.wikipedia.org/wiki/Hydraulic%20machinery | Hydraulic machinery | Hydraulic machines use liquid fluid power to perform work. Heavy construction vehicles are a common example. In this type of machine, hydraulic fluid is pumped to various hydraulic motors and hydraulic cylinders throughout the machine and becomes pressurized according to the resistance present. The fluid is controlled directly or automatically by control valves and distributed through hoses, tubes, or pipes.
Hydraulic systems, like pneumatic systems, are based on Pascal's law which states that any pressure applied to a fluid inside a closed system will transmit that pressure equally everywhere and in all directions. A hydraulic system uses an incompressible liquid as its fluid, rather than a compressible gas.
The popularity of hydraulic machinery is due to the large amount of power that can be transferred through small tubes and flexible hoses, the high power density and a wide array of actuators that can make use of this power, and the huge multiplication of forces that can be achieved by applying pressures over relatively large areas. One drawback, compared to machines using gears and shafts, is that any transmission of power results in some losses due to resistance of fluid flow through the piping.
History
Joseph Bramah patented the hydraulic press in 1795. While working at Bramah's shop, Henry Maudslay suggested a cup leather packing. Because it produced superior results, the hydraulic press eventually displaced the steam hammer for metal forging.
To supply large-scale power that was impractical for individual steam engines, central station hydraulic systems were developed. Hydraulic power was used to operate cranes and other machinery in British ports and elsewhere in Europe. The largest hydraulic system was in London. Hydraulic power was used extensively in Bessemer steel production. Hydraulic power was also used for elevators, to operate canal locks and rotating sections of bridges. Some of these systems remained in use well into the twentieth century.
Harry Franklin Vickers was called the "Father of Industrial Hydraulics" by ASME.
Force and torque multiplication
A fundamental feature of hydraulic systems is the ability to apply force or torque multiplication in an easy way, independent of the distance between the input and output, without the need for mechanical gears or levers, either by altering the effective areas in two connected cylinders or the effective displacement (cc/rev) between a pump and motor. In normal cases, hydraulic ratios are combined with a mechanical force or torque ratio for optimum machine designs such as boom movements and track drives for an excavator.
Examples
Two hydraulic cylinders interconnected
Cylinder C1 is one inch in radius, and cylinder C2 is ten inches in radius. If the force exerted on C1 is 10 lbf, the force exerted by C2 is 1000 lbf because C2 is a hundred times larger in area (S = πr²) as C1. The downside to this is that you have to move C1 a hundred inches to move C2 one inch. The most common use for this is the classical hydraulic jack where a pumping cylinder with a small diameter is connected to the lifting cylinder with a large diameter.
Pump and motor
If a hydraulic rotary pump with the displacement 10 cc/rev is connected to a hydraulic rotary motor with 100 cc/rev, the shaft torque required to drive the pump is one-tenth of the torque then available at the motor shaft, but the shaft speed (rev/min) for the motor is also only one-tenth of the pump shaft speed. This combination is actually the same type of force multiplication as the cylinder example, just that the linear force in this case is a rotary force, defined as torque.
Both these examples are usually referred to as a hydraulic transmission or hydrostatic transmission involving a certain hydraulic "gear ratio".
Hydraulic circuits
A hydraulic circuit is a system comprising an interconnected set of discrete components that transport liquid. The purpose of this system may be to control where fluid flows (as in a network of tubes of coolant in a thermodynamic system) or to control fluid pressure (as in hydraulic amplifiers). For example, hydraulic machinery uses hydraulic circuits (in which hydraulic fluid is pushed, under pressure, through hydraulic pumps, pipes, tubes, hoses, hydraulic motors, hydraulic cylinders, and so on) to move heavy loads. The approach of describing a fluid system in terms of discrete components is inspired by the success of electrical circuit theory. Just as electric circuit theory works when elements are discrete and linear, hydraulic circuit theory works best when the elements (passive components such as pipes or transmission lines or active components such as power packs or pumps) are discrete and linear. This usually means that hydraulic circuit analysis works best for long, thin tubes with discrete pumps, as found in chemical process flow systems or microscale devices.
The circuit comprises the following components:
Active components
Hydraulic power pack
Transmission lines
Hydraulic hoses
Passive components
Hydraulic cylinders
For the hydraulic fluid to do work, it must flow to the actuator and/or motors, then return to a reservoir. The fluid is then filtered and re-pumped. The path taken by hydraulic fluid is called a hydraulic circuit of which there are several types.
Open center circuits use pumps that supply a continuous flow. The flow is returned to tank through the control valve's open center; that is, when the control valve is centered, it provides an open return path to the tank and the fluid is not pumped to high pressure. Otherwise, if the control valve is actuated it routes fluid to and from an actuator and tank. The fluid's pressure will rise to meet any resistance, since the pump has a constant output. If the pressure rises too high, fluid returns to the tank through a pressure relief valve. Multiple control valves may be stacked in series. This type of circuit can use inexpensive, constant displacement pumps.
Closed center circuits supply full pressure to the control valves, whether any valves are actuated or not. The pumps vary their flow rate, pumping very little hydraulic fluid until the operator actuates a valve. The valve's spool therefore doesn't need an open center return path to tank. Multiple valves can be connected in a parallel arrangement and system pressure is equal for all valves.
Open loop circuits
Open-loop: Pump-inlet and motor-return (via the directional valve) are connected to the hydraulic tank. The term loop applies to feedback; the more correct term is open versus closed "circuit". Open center circuits use pumps which supply a continuous flow. The flow is returned to the tank through the control valve's open center; that is, when the control valve is centered, it provides an open return path to the tank and the fluid is not pumped to a high pressure. Otherwise, if the control valve is actuated it routes fluid to and from an actuator and tank. The fluid's pressure will rise to meet any resistance, since the pump has a constant output. If the pressure rises too high, fluid returns to the tank through a pressure relief valve. Multiple control valves may be stacked in series. This type of circuit can use inexpensive, constant displacement pumps.
Closed loop circuits
Closed-loop: Motor-return is connected directly to the pump-inlet. To keep up pressure on the low pressure side, the circuits have a charge pump (a small gear pump) that supplies cooled and filtered oil to the low pressure side. Closed-loop circuits are generally used for hydrostatic transmissions in mobile applications. Advantages: No directional valve and better response, the circuit can work with higher pressure. The pump swivel angle covers both positive and negative flow direction. Disadvantages: The pump cannot be utilized for any other hydraulic function in an easy way and cooling can be a problem due to limited exchange of oil flow. High power closed loop systems generally must have a 'flush-valve' assembled in the circuit in order to exchange much more flow than the basic leakage flow from the pump and the motor, for increased cooling and filtering. The flush valve is normally integrated in the motor housing to get a cooling effect for the oil that is rotating in the motor housing itself. The losses in the motor housing from rotating effects and losses in the ball bearings can be considerable as motor speeds will reach 4000-5000 rev/min or even more at maximum vehicle speed. The leakage flow as well as the extra flush flow must be supplied by the charge pump. A large charge pump is thus very important if the transmission is designed for high pressures and high motor speeds. High oil temperature is usually a major problem when using hydrostatic transmissions at high vehicle speeds for longer periods, for instance when transporting the machine from one work place to the other. High oil temperatures for long periods will drastically reduce the lifetime of the transmission. To keep down the oil temperature, the system pressure during transport must be lowered, meaning that the minimum displacement for the motor must be limited to a reasonable value. Circuit pressure during transport around 200-250 bar is recommended.
Closed loop systems in mobile equipment are generally used for the transmission as an alternative to mechanical and hydrodynamic (converter) transmissions. The advantage is a stepless gear ratio (continuously variable speed/torque) and a more flexible control of the gear ratio depending on the load and operating conditions. The hydrostatic transmission is generally limited to around 200 kW maximum power, as the total cost gets too high at higher power compared to a hydrodynamic transmission. Large wheel loaders for instance and heavy machines are therefore usually equipped with converter transmissions. Recent technical achievements for the converter transmissions have improved the efficiency and developments in the software have also improved the characteristics, for example selectable gear shifting programs during operation and more gear steps, giving them characteristics close to the hydrostatic transmission.
Constant pressure and load-sensing systems
Hydrostatic transmissions for earth moving machines, such as for track loaders, are often equipped with a separate 'inch pedal' that is used to temporarily increase the diesel engine rpm while reducing the vehicle speed in order to increase the available hydraulic power output for the working hydraulics at low speeds and increase the tractive effort. The function is similar to stalling a converter gearbox at high engine rpm. The inch function affects the preset characteristics for the 'hydrostatic' gear ratio versus diesel engine rpm.
Constant pressure systems
The closed center circuits exist in two basic configurations, normally related to the regulator for the variable pump that supplies the oil:
Constant pressure systems (CP), standard. Pump pressure always equals the pressure setting for the pump regulator. This setting must cover the maximum required load pressure. Pump delivers flow according to required sum of flow to the consumers. The CP system generates large power losses if the machine works with large variations in load pressure and the average system pressure is much lower than the pressure setting for the pump regulator. CP is simple in design, and works like a pneumatic system. New hydraulic functions can easily be added and the system is quick in response.
Constant pressure systems, unloaded. Same basic configuration as 'standard' CP system but the pump is unloaded to a low stand-by pressure when all valves are in neutral position. Not so fast response as standard CP but pump lifetime is prolonged.
Load-sensing systems
Load-sensing systems (LS) generate less power losses as the pump can reduce both flow and pressure to match the load requirements, but require more tuning than the CP system with respect to system stability. The LS system also requires additional logical valves and compensator valves in the directional valves, thus it is technically more complex and more expensive than the CP system. The LS system generates a constant power loss related to the regulating pressure drop for the pump regulator :
The average is around 2 MPa (290 psi). If the pump flow is high the extra loss can be considerable. The power loss also increases if the load pressures vary a lot. The cylinder areas, motor displacements and mechanical torque arms must be designed to match load pressure in order to bring down the power losses. Pump pressure always equals the maximum load pressure when several functions are run simultaneously and the power input to the pump equals the (max. load pressure + ΔpLS) x sum of flow.
Five basic types of load sensing systems
Load sensing without compensators in the directional valves. Hydraulically controlled LS pump.
Load sensing with up-stream compensator for each connected directional valve. Hydraulically controlled LS pump.
Load sensing with down-stream compensator for each connected directional valve. Hydraulically controlled LS pump.
Load sensing with a combination of up-stream and down-stream compensators. Hydraulically controlled LS pump.
Load sensing with synchronized, both electric controlled pump displacement and electric controlled valve flow area for faster response, increased stability and fewer system losses. This is a new type of LS-system, not yet fully developed.
Technically the down-stream mounted compensator in a valve block can physically be mounted "up-stream", but work as a down-stream compensator.
System type (3) gives the advantage that activated functions are synchronized independent of pump flow capacity. The flow relation between two or more activated functions remains independent of load pressures, even if the pump reaches the maximum swivel angle. This feature is important for machines that often run with the pump at maximum swivel angle and with several activated functions that must be synchronized in speed, such as with excavators. With the type (4) system, the functions with up-stream compensators have priority, for example the steering function for a wheel loader. The system type with down-stream compensators usually have a unique trademark depending on the manufacturer of the valves, for example "LSC" (Linde Hydraulics), "LUDV" (Bosch Rexroth Hydraulics) and "Flowsharing" (Parker Hydraulics) etc. No official standardized name for this type of system has been established but flowsharing is a common name for it.
Components
Hydraulic pump
Hydraulic pumps supply fluid to the components in the system. Pressure in the system develops in reaction to the load. Hence, a pump rated for 5,000 psi is capable of maintaining flow against a load of 5,000 psi.
Pumps have a power density about ten times greater than an electric motor (by volume). They are powered by an electric motor or an engine, connected through gears, belts, or a flexible elastomeric coupling to reduce vibration.
Common types of hydraulic pumps to hydraulic machinery applications are:
Gear pump: cheap, durable (especially in g-rotor form), simple. Less efficient, because they are constant (fixed) displacement, and mainly suitable for pressures below 20 MPa (3000 psi).
Vane pump: cheap and simple, reliable. Good for higher-flow low-pressure output.
Axial piston pump: many designed with a variable displacement mechanism, to vary output flow for automatic control of pressure. There are various axial piston pump designs, including swashplate (sometimes referred to as a valveplate pump) and checkball (sometimes referred to as a wobble plate pump). The most common is the swashplate pump. A variable-angle swashplate causes the pistons to reciprocate a greater or lesser distance per rotation, allowing output flow rate and pressure to be varied (greater displacement angle causes higher flow rate, lower pressure, and vice versa).
Radial piston pump: normally used for very high pressure at small flows.
Piston pumps are more expensive than gear or vane pumps, but provide longer life operating at higher pressure, with difficult fluids and longer continuous duty cycles. Piston pumps make up one half of a hydrostatic transmission.
Control valves
Directional control valves route the fluid to the desired actuator. They usually consist of a spool inside a cast iron or steel housing. The spool slides to different positions in the housing, and intersecting grooves and channels route the fluid based on the spool's position.
The spool has a central (neutral) position maintained with springs; in this position the supply fluid is blocked, or returned to tank. Sliding the spool to one side routes the hydraulic fluid to an actuator and provides a return path from the actuator to tank. When the spool is moved to the opposite direction the supply and return paths are switched. When the spool is allowed to return to neutral (center) position the actuator fluid paths are blocked, locking it in position.
Directional control valves are usually designed to be stackable, with one valve for each hydraulic cylinder, and one fluid input supplying all the valves in the stack.
Tolerances are very tight in order to handle the high pressure and avoid leaking, spools typically have a clearance with the housing of less than a thousandth of an inch (25 μm). The valve block will be mounted to the machine's frame with a three point pattern to avoid distorting the valve block and jamming the valve's sensitive components.
The spool position may be actuated by mechanical levers, hydraulic pilot pressure, or solenoids which push the spool left or right. A seal allows part of the spool to protrude outside the housing, where it is accessible to the actuator.
The main valve block is usually a stack of off the shelf directional control valves chosen by flow capacity and performance. Some valves are designed to be proportional (flow rate proportional to valve position), while others may be simply on-off. The control valve is one of the most expensive and sensitive parts of a hydraulic circuit.
Pressure relief valves are used in several places in hydraulic machinery; on the return circuit to maintain a small amount of pressure for brakes, pilot lines, etc... On hydraulic cylinders, to prevent overloading and hydraulic line/seal rupture. On the hydraulic reservoir, to maintain a small positive pressure which excludes moisture and contamination.
Pressure regulators reduce the supply pressure of hydraulic fluids as needed for various circuits.
Sequence valves control the sequence of hydraulic circuits; to ensure that one hydraulic cylinder is fully extended before another starts its stroke, for example. Hydraulic circuits can perform a sequence of operations automatically, such as trip-and-reclose three times, then lockout, of an oil-interrupting recloser.
Shuttle valves provide a logical or function.
Check valves are one-way valves, allowing an accumulator to charge and maintain its pressure after the machine is turned off, for example.
Pilot controlled check valves are one-way valve that can be opened (for both directions) by a foreign pressure signal. For instance if the load should not be held by the check valve anymore. Often the foreign pressure comes from the other pipe that is connected to the motor or cylinder.
Counterbalance valves are in fact a special type of pilot controlled check valve. Whereas the check valve is open or closed, the counterbalance valve acts a bit like a pilot controlled flow control.
Cartridge valves are in fact the inner part of a check valve; they are off the shelf components with a standardized envelope, making them easy to populate a proprietary valve block. They are available in many configurations; on/off, proportional, pressure relief, etc. They generally screw into a valve block and are electrically controlled to provide logic and automated functions.
Hydraulic fuses are in-line safety devices designed to automatically seal off a hydraulic line if pressure becomes too low, or safely vent fluid if pressure becomes too high.
Auxiliary valves in complex hydraulic systems may have auxiliary valve blocks to handle various duties unseen to the operator, such as accumulator charging, cooling fan operation, air conditioning power, etc. They are usually custom valves designed for the particular machine, and may consist of a metal block with ports and channels drilled. Cartridge valves are threaded into the ports and may be electrically controlled by switches or a microprocessor to route fluid power as needed.
Actuators
Hydraulic cylinder
Hydraulic motor (a pump plumbed in reverse); hydraulic motors with axial configuration use swashplates for highly accurate control and also in 'no stop' continuous (360°) precision positioning mechanisms. These are frequently driven by several hydraulic pistons acting in sequence.
Hydrostatic transmission
Brakes
Reservoir
The hydraulic fluid reservoir holds excess hydraulic fluid to accommodate volume changes from: cylinder extension and contraction, temperature driven expansion and contraction, and leaks. The reservoir is also designed to aid in separation of air from the fluid and also work as a heat accumulator to cover losses in the system when peak power is used. Reservoirs can also help separate dirt and other particulate from the oil, as the particulate will generally settle to the bottom of the tank.
Some designs include dynamic flow channels on the fluid's return path that allow for a smaller reservoir.
Accumulators
Accumulators are a common part of hydraulic machinery. Their function is to store energy by using pressurized gas. One type is a tube with a floating piston. On the one side of the piston there is a charge of pressurized gas, and on the other side is the fluid. Bladders are used in other designs. Reservoirs store a system's fluid.
Examples of accumulator uses are backup power for steering or brakes, or to act as a shock absorber for the hydraulic circuit.
Hydraulic fluid
Also known as tractor fluid, hydraulic fluid is the life of the hydraulic circuit. It is usually petroleum oil with various additives. Some hydraulic machines require fire resistant fluids, depending on their applications. In some factories where food is prepared, either an edible oil or water is used as a working fluid for health and safety reasons.
In addition to transferring energy, hydraulic fluid needs to lubricate components, suspend contaminants and metal filings for transport to the filter, and to function well to several hundred degrees Fahrenheit or Celsius.
Filters
Filters are an important part of hydraulic systems which removes the unwanted particles from fluid. Metal particles are continually produced by mechanical components and need to be removed along with other contaminants.
Filters may be positioned in many locations. The filter may be located between the reservoir and the pump intake. Blockage of the filter will cause cavitation and possibly failure of the pump. Sometimes the filter is located between the pump and the control valves. This arrangement is more expensive, since the filter housing is pressurized, but eliminates cavitation problems and protects the control valve from pump failures. The third common filter location is just before the return line enters the reservoir. This location is relatively insensitive to blockage and does not require a pressurized housing, but contaminants that enter the reservoir from external sources are not filtered until passing through the system at least once. Filters are used from 7 micron to 15 micron depends upon the viscosity grade of hydraulic oil.
Tubes, pipes and hoses
Hydraulic tubes are seamless steel precision pipes, specially manufactured for hydraulics. The tubes have standard sizes for different pressure ranges, with standard diameters up to 100 mm. The tubes are supplied by manufacturers in lengths of 6 m, cleaned, oiled and plugged. The tubes are interconnected by different types of flanges (especially for the larger sizes and pressures), welding cones/nipples (with o-ring seal), several types of flare connection and by cut-rings. In larger sizes, hydraulic pipes are used. Direct joining of tubes by welding is not acceptable since the interior cannot be inspected.
Hydraulic pipe is used in case standard hydraulic tubes are not available. Generally these are used for low pressure. They can be connected by threaded connections, but usually by welds. Because of the larger diameters the pipe can usually be inspected internally after welding. Black pipe is non-galvanized and suitable for welding.
Hydraulic hose is graded by pressure, temperature, and fluid compatibility. Hoses are used when pipes or tubes can not be used, usually to provide flexibility for machine operation or maintenance. The hose is built up with rubber and steel layers. A rubber interior is surrounded by multiple layers of woven wire and rubber. The exterior is designed for abrasion resistance. The bend radius of hydraulic hose is carefully designed into the machine, since hose failures can be deadly, and violating the hose's minimum bend radius will cause failure. Hydraulic hoses generally have steel fittings swaged on the ends. The weakest part of the high pressure hose is the connection of the hose to the fitting. Another disadvantage of hoses is the shorter life of rubber which requires periodic replacement, usually at five to seven year intervals.
Tubes and pipes for hydraulic n applications are internally oiled before the system is commissioned. Usually steel piping is painted outside. Where flare and other couplings are used, the paint is removed under the nut, and is a location where corrosion can begin. For this reason, in marine applications most piping is stainless steel.
Seals, fittings and connections
Components of a hydraulic system [sources (e.g. pumps), controls (e.g. valves) and actuators (e.g. cylinders)] need connections that will contain and direct the hydraulic fluid without leaking or losing the pressure that makes them work. In some cases, the components can be made to bolt together with fluid paths built-in. In more cases, though, rigid tubing or flexible hoses are used to direct the flow from one component to the next. Each component has entry and exit points for the fluid involved (called ports) sized according to how much fluid is expected to pass through it.
There are a number of standardized methods in use to attach the hose or tube to the component. Some are intended for ease of use and service, others are better for higher system pressures or control of leakage. The most common method, in general, is to provide in each component a female-threaded port, on each hose or tube a female-threaded captive nut, and use a separate adapter fitting with matching male threads to connect the two. This is functional, economical to manufacture, and easy to service.
Fittings serve several purposes;
To join components with ports of different sizes.
To bridge different standards; O-ring boss to JIC, or pipe threads to face seal, for example.
To allow proper orientation of components, a 90°, 45°, straight, or swivel fitting is chosen as needed. They are designed to be positioned in the correct orientation and then tightened.
To incorporate bulkhead hardware to pass the fluid through an obstructing wall.
A quick disconnect fitting may be added to a machine without modification of hoses or valves
A typical piece of machinery or heavy equipment may have thousands of sealed connection points and several different types:
Pipe fittings, the fitting is screwed in until tight, difficult to orient an angled fitting correctly without over or under tightening.
O-ring boss, the fitting is screwed into a boss and orientated as needed, an additional nut tightens the fitting, washer and o-ring in place.
Flare fittings, are metal to metal compression seals deformed with a cone nut and pressed into a flare mating.
Face seal, metal flanges with a groove and o-ring seal are fastened together.
Beam seals are costly metal to metal seals used primarily in aircraft.
Swaged seals, tubes are connected with fittings that are swaged permanently in place. Primarily used in aircraft.
Elastomeric seals (O-ring boss and face seal) are the most common types of seals in heavy equipment and are capable of reliably sealing more than of fluid pressure.
| Technology | Hydraulics and pneumatics | null |
1372941 | https://en.wikipedia.org/wiki/Shenzhen%20Metro | Shenzhen Metro | The Shenzhen Metro () is the rapid transit system for the city of Shenzhen in Guangdong province, China. The newest lines and extensions which opened on December 27, 2024 put the network at of trackage. It currently operates on 17 lines with 398 stations. Despite having only opened on December 28, 2004, the Shenzhen Metro is the 5th longest metro system in the world. By 2035, the network is planned to comprise 8 express and 24 non-express lines totaling of trackage.
Current system
Currently the network has of route, operating on 17 lines with 398 stations. Line 1, Line 4 and Line 10 run to the border crossings between the Shenzhen Special Economic Zone and the Hong Kong Special Administrative Region at Luohu/Lo Wu and Futian Checkpoint/Lok Ma Chau, where riders can transfer to Hong Kong's MTR East Rail line for travel onwards to Hong Kong.
Line 1
Line 1, formerly known as Luobao line, runs westward from Luohu to Airport East. Trains operate every 2 minutes during peak hours and every 4 minutes at other times. The line is operated by SZMC (Shenzhen Metro Group). Line 1's color is green.
December 28, 2004: Luohu – Window of the World
September 28, 2009: Window of the World – Shenzhen University
June 15, 2011: Shenzhen University – Airport East
Line 2
Line 2, formerly known as Shekou line, runs from Chiwan to Liantang. Line 2 is connected with Line 8 at Liantang station. The line is operated by SZMC (Shenzhen Metro Group). Line 2's color is dark orange, the same as Line 8.
December 28, 2010: Chiwan – Window of the World
June 28, 2011: Window of the World – Xinxiu
October 28, 2020: Xinxiu – Liantang
Line 3
Line 3, formerly known as Longgang line, runs from Futian Bonded Area to Shuanlong in Longgang, in the north-east part of the city. Construction began on December 26, 2005. The line is operated by Shenzhen Metro Line 3 Operations, which has been a subsidiary of SZMC (Shenzhen Metro Group) since April 11, 2011, when an 80% stake was transferred to SZMC. Line 3's color is sky blue.
December 28, 2010: Caopu – Shuanglong
June 28, 2011: Yitian – Caopu
October 28, 2020: Futian Bonded Area – Yitian
Line 4
Line 4, formerly known as Longhua line, runs northward from Futian Checkpoint to Niuhu. Trains operate every 2.5 minutes at peak hours and every 6 minutes during off-peak hours. Stations from Futian Checkpoint to Shangmeilin Station are underground. The line has been operated by MTR Corporation (Shenzhen), a subsidiary of MTR Corporation, since July 1, 2010. Line 4's color is red.
December 28, 2004: Fumin – Children's Palace
June 28, 2007: Futian Checkpoint – Fumin
June 16, 2011: Children's Palace – Qinghu
October 28, 2020: Qinghu – Niuhu
Line 5
Line 5, formerly known as Huanzhong line, runs from Chiwan in the west to Huangbeiling in the east. Construction began in May 2009 and the line opened on June 22, 2011. Line 5 required a total investment of 20.6 billion RMB. The line is operated by SZMC (Shenzhen Metro Group). Line 5's color is purple.
June 22, 2011: Qianhaiwan – Huangbeiling
September 28, 2019: Qianhaiwan – Chiwan
Line 6
Line 6, formerly known as Guangming line, runs from Songgang in the north to Science Museum in the south, with a length of and a total of 27 stations. Construction began in August 2015 and the line opened on August 18, 2020. The line is operated by SZMC (Shenzhen Metro Group). Line 6's color is mint green.
August 18, 2020: Science Museum – Songgang
Line 6 Branch
Line 6 Branch, also known as Branch Line 6, runs from Guangming to SIAT in the north. The line opened on November 28, 2022. Line 6 Branch's color is teal.
November 28, 2022: Guangming – SIAT
Line 7
Line 7, formerly known as Xili line of the Shenzhen Metro, opened on October 28, 2016, with a length of and a total of 27 stations. It connects the Xili Lake to Tai'an. The line travels East–West across Shenzhen in a "V" shape. The line is operated by SZMC (Shenzhen Metro Group). Line 7's color is navy blue.
October 28, 2016: Xili Lake – Tai'an
Line 8
Line 8, formerly known as Yantian line of the Shenzhen Metro, opened on October 28, 2020, with a length of and a total of 11 stations. It connects the eastern suburbs of Liantang to Yantian Road, then towards the beach resorts at Dameisha and Xiaomeisha. However, this line serves as the extension of Line 2 in actual operation. The line is operated by SZMC (Shenzhen Metro Group). Line 8's color is dark orange, the same as Line 2.
October 28, 2020: Liantang – Yantian Road
December 27, 2023: Yantian Road – Xiaomeisha
Line 9
Line 9, formerly known as Meilin line or Neihuan line of the Shenzhen Metro, opened on October 28, 2016. The line runs eastward from to . It has 10 transfer stations. The line is long, running through the districts of Nanshan, Futian and Luohu. The line is operated by SZMC (Shenzhen Metro Group). Line 9's color is grey brown.
October 28, 2016: –
December 8, 2019: –
Line 10
Line 10 formerly known as Bantian line, runs from Futian Checkpoint in the south to Shuangyong Street in the north, with a length of and a total of 24 stations. Construction began in September 2015 and the line opened on August 18, 2020. The line is operated by SZMC (Shenzhen Metro Group). Line 10's color is pink.
August 18, 2020: Futian Checkpoint – Shuangyong Street
Line 11
Line 11, also known as the Airport Express, runs from Bitou in the northwest to Gangxia North in the city centre via Shenzhen Bao'an International Airport. Construction began in April 2012 and the line opened on June 28, 2016. Line 11 runs at a higher speed of , however, the current line speed is limited to . The line is operated by SZMC (Shenzhen Metro Group). Line 11's color is maroon.
June 28, 2016: Bitou – Futian
October 28, 2022: Futian – Gangxia North
December 28, 2024: Gangxia North – Huaqiang South
Line 12
Line 12, also known as Nanbao Line, runs from Zuopaotai East in the southwest to Waterlands Resort East in the northwest. Construction began in 2018 and the line opened on November 28, 2022. Line 12's color is light purple.
November 28, 2022: Zuopaotai East – Waterlands Resort East
Line 13
Line 13, also known as Shiyan Line, runs from Shenzhen Bay Checkpoint at the Shenzhen Bay Port in Nanshan to Shangwu in northeast Bao'an. Construction began in 2018 and the first phase of the line between Shenzhen Bay Checkpoint and Hi-Tech Central opened on December 28, 2024. Line 13's color is light orange.
December 28, 2024: –
Line 14
Line 14, also known as the Eastern Express, runs from Gangxia North in the city centre to Shatian in the northeast. Construction began in 2018 and the line opened on October 28, 2022. The line is operated by SZMC (Shenzhen Metro Group). Line 14's color is yellow.
October 28, 2022: Gangxia North – Shatian
Line 16
Line 16, also known as Longping Line, runs from Universiade in the centre of Longgang to Tianxin in the northeast. Construction began in 2018 and the line opened on December 28, 2022. The line is operated by SZMC (Shenzhen Metro Group). Line 16's color is dark blue.
December 28, 2022: Universiade – Tianxin
Line 20
Line 20, formerly known as Fuyong line, runs from Airport North in the north-west to Convention & Exhibition City near Shenzhen World. Construction began in September 2016 and the line opened on December 28, 2021. Line 20 runs at the same top speed as line 11, at . The line is operated by SZMC (Shenzhen Metro Group). Line 20's color is light blue.
December 28, 2021: Airport North – Convention & Exhibition City
History
Early planning
In late 1983, Party Secretary of Shenzhen Mayor Liang Xiang led a team to Singapore to study its mass transit system. Upon returning it was decided that on each side of Shennan Avenue should be protected as a green belt, and to set aside a wide median reserved for a light rail or light metro line. In 1984, the "Shenzhen Special Economic Zone Master Plan (1985–2000)" pointed out that, with the growing population and traffic in Shenzhen, a light metro system would not have sufficient capacity to meet future demand. Instead the report proposed a heavy rail subway line to be built along Shennan Avenue. The project was finally approved by the Central Planning Department in 1992.
In August 1992, during and re-feasibility and rail network planning, The Shenzhen Municipal Government decided to move from building a light metro line to a heavy rail subway line. The rapid growth of Shenzhen City made a lower capacity light metro line impractical. In 1994, Shenzhen organized the preparation of the "Shenzhen urban rail network master plan" to be incorporated into the "Shenzhen City Master Plan (1996–2010)". The city's vision for an urban rail network would consists of nine lines. Of the nine transit lines, three of them would be commuter rail lines upgraded from existing national mainline railways. The total length of the proposed network would be about . The three upgraded commuter rail lines would overlap the Guangzhou–Shenzhen railway, Pinghu–Nanshan railway and Pingyan railway. This plan established the basic framework for the Shenzhen Metro network.
Construction suspended and restarted
In December 1995, the State Council issued the "moratorium on approval of urban rapid transit projects" to suspend approval of rail transit projects in all Chinese cities except Beijing, Shanghai, and Guangzhou. The Shenzhen Metro project was postponed. In 1996, prior to the handover of Hong Kong, authorities attempted to restart construction by renaming the project "The Luohu, Huanggang / Lok Ma Chau border crossing passenger rail connection project", stressing that the project is designed to meet the potential growing demand for cross-border passenger traffic after the handover.
In 1997, Shenzhen reapplied its Subway plans to the State Planning Commission, and received approval in May 1998. The project was renamed the "Shenzhen Metro first phase". In July 1998, SZMC (Shenzhen Metro Group). was formally established. By April 1999, the subway project feasibility study report has been approved by the state.
Phase I (1998–2004)
Construction of the first sections of Line 1 and Line 4 began in 1999. The grand opening of the Shenzhen Metro system occurred at 5:00 pm on Tuesday, December 28, 2004. This made Shenzhen the seventh city in mainland China to have a subway after Beijing, Tianjin, Shanghai, Guangzhou, Dalian and Wuhan.
Initially the trains operated at 15-minute frequencies and consisted of Line 1 services between Luohu and Shijie Zhi Chuang (now Window of the World) and the Line 4 services between Fumin and Shaonian Gong (now Children's Palace). Initially the English names of the stations were rendered in Hanyu Pinyin, but some of the names were changed to English translation with American spelling like the rest of mainland China, despite being close to the Hong Kong, which uses British spelling and ongoing political tensions with the US, in mid-2011.
The Futian Checkpoint station opened on June 28, 2007, using the name Huanggang.
Name changes
On April 23, 2008, Shenzhen Municipal Planning Bureau announced that it would change the nomenclature of Shenzhen's subway lines according to the "2007 Urban Rail Transit Plan Scheme". Instead of using numbers as the lines official designation, as typically used in other mainland Chinese metro systems, lines would be given Chinese names more akin to the Hong Kong MTR. In 2010, the Scheme was reviewed and adjusted with new routes and names in addition to newly proposed lines. On October 23, 2013, the SZMC (Shenzhen Metro Group) decided that current operational lines will have their number and names combined, while future lines will only be numbered. Due to the change in the construction order of several lines, some numerical names have been reviewed in order to prevent big jump between numbers. By 2016, only numerical names are used.
Lines currently in operation:
Lines under construction:
Phase II (2007–2011)
From 2004 to 2007, there was a lack of official government interest and attention to expanding the subway after completion of Phase 1 with little or no active projects. Subway construction speed was ridiculed as "earthworm speed". On January 17, 2007, Shenzhen won the right to host the 2011 Universiade. In the bid Shenzhen committed to complete of subway lines before the games. The mayor of Shenzhen at the time, Xu Zongheng, sharply criticized the speed and efficiency of Shenzhen's subway construction procedures and calls for reform. Subsequently, the Shenzhen municipal government and various departments signed a liability form, requiring Phase II subway expansion to be completed in time for the Universiade. Shenzhen Metro increased to over a hundred operating metro stations in June 2011, just before the Shenzhen Universiade games. In the span of two weeks, the network expanded from to . This expansion increased rail transit's share of total public transit trips from 6% to 29% in 2014.
Phase III (2012–2023)
In 2010, the Shenzhen Urban Planning and Land Resources Committee proposed a building program (Phase III) between 2011 and 2020. In 2011 this plan was approved by the NDRC. Phase III formally commenced in May 2011 with an expected cost of 125.6 billion yuan. It will cover Lines 6, 7, 8, 9, and 11 and will extend the length of the Shenzhen Metro to and 10 lines. In June 2011, the Shenzhen Urban Planning and Land Resources Commission started gather public input on Phase III station names. On June 28, 2016, Line 11 opened being the first subway line in Shenzhen with 8 car trains and maximum service speed and the first in China with a First Class service. Lines 7 and 9 followed on October 28, 2016. South extension of Line 5 opened on September 28, 2019, and west extension of Line 9 opened on December 8, 2019. Line 6 and Line 10 opened on August 18, 2020, bringing the length of the Shenzhen Metro to and the fourth longest in China. Second east extension of Line 2, second south extension of Line 3, second north extension of Line 4 and phase 1 of Line 8 opened on October 28, 2020, bringing the length of the Shenzhen Metro to . Phase III is also the first phase in which the lines are officially numbered instead of named and colored.
Phase IV & Phase IV revised expansion (2017–2026)
With the shortening of the Phase III implementation period, a number of lines (Lines 16 and 12) planned in 2007's Phase III moved into the next phase. By 2016, it was determined that Phase IV will have an implementation period between 2017 and 2022 and consist of of new subway. Lines 13 and 14 which originally had a long term 2030 completion deadline were moved into Phase IV expansion. In addition, a branch line of Line 6 will connect with the neighboring Dongguan Rail Transit system. Lines 12, 13, 14, and 16 and branch of Line 6 was approved by the NDRC in July 2017 and started construction in January 2018. The first phase of Line 20 was fast tracked from Phase IV to provide a shuttle between Line 11 and a new International Convention Center, now called Convention & Exhibition City. The construction started in September 2016, but as for early 2019, the construction is paused because the Development and Reform Commission did not approve the project. The Phase IV revised plan approved by the NDRC on March 26, 2020, approved the first Phase of Line 20 allowing for construction to continue. The line eventually opened on December 28, 2021. Futian to Gangxia North in the first section of Phase 2 of Line 11 opened on October 28, 2022 in tandem with Line 14 Phase I.
Future expansion
Phase IV expansion
Phase IV Revised Expansion
The Phase IV revised plan approved by the NDRC on March 26, 2020, added a number of extension projects.
Miscellaneous
Line 5 west extension is part of Phase II expansion.
Phase V (2023–2028)
In the Shenzhen Metro 2007 masterplan proposed four more lines (Lines 13, 14, 15 and 16) which have a planned completion target of 2030. In 2016, all aforementioned lines but Line 15 were designated as part of the Phase IV expansion, moving the completion date forward from 2030 to 2022. In 2012, four further lines Qiannan (Line 17), Pinghu (Line 18), Pingshan (Line 19) and Fuyong (Line 20) where unveiled, making the total planned length of the Shenzhen Metro to spread out over 20 lines. The first phase of Line 20 was fast-tracked and included in the Phase III revised expansion with a completion date of 2018. This leaves Line 15, 17–19 and the rest of Line 20 available for the next phase (Phase V) of subway expansion. In September 2022, the Shenzhen municipal government confirmed the projects proposed to be included in its phase V expansion. A total of of new lines are proposed.
On March 31, 2023, the Bureau of Housing and Urban Rural Development of Shenzhen Metro Municipality will open the bidding for the fifth phase planning of Shenzhen Metro Metro, including Line 15, Line 17 Phase I, Line 19 Phase I, Line 20 Phase II (Airport East—Baishizhou), Line 22 Phase I, Line 25 Phase I, Line 27 Phase I, Line 29 Phase I and Line 32 Phase I, This means that these 9 lines have been approved with an investment amount of 191.1 billion yuan. However, Line 18 Phase I, Line 21 Phase I, Line 10 East Extension (Shenzhen Section) and Metro Line 11 North Extension (Shenzhen Section), which were previously proposed to be included in the fifth phase plan of Shenzhen Metro, were not included in this announcement.
In June 2023, the Fifth Phase Construction Plan of Shenzhen Urban Rail Transit (2023–2028) has been approved for a total of 11 construction projects, including Line 15, Line 17 Phase I, Line 19 Phase I, Line 20 Phase II (Airport East—Baishizhou), Line 22 Phase I, Line 25 Phase I, Line 27 Phase I, Line 29 Phase I, Line 32 Phase I, Line 10 East Extension (Shenzhen Section) and Metro Line 11 North Extension (Shenzhen Section).
Long-term plan
Aside from the set masterplan, at the 12th Guangdong Provincial People's Congress in January 2014, it was proposed to extend Line 4 beyond the planned Phase III terminus at the Songyuan Bus Terminal in Guanlan. The proposal wanted to further extend this line to reach the future planned Dongguan Metro Line 4 at Tangxia station. This proposal aims to shorten the distance between the two cities in residents' minds, boost tourism industries in both cities and expand housing options. It would also allow for direct connection between Hong Kong and Dongguan. As the area in the proposed area is less developed, the cost in building the line is expected to be lower, with a feasibility study yet to be conducted. In addition to metro lines, 5 Pearl River Delta Rapid Transit lines connecting neighboring urban centers in the Pearl River Delta such as Dongguan, Huizhou, Foshan and Guangzhou, totaling , have also been revealed. In 2016, an even more ambitious masterplan, expanding the previously planned 20 lines to 32, was unveiled. The new plan envisions a subway network to be completed by 2030. This will allow for travel between the central and suburban districts to be shortened to 45 minutes and for public transit to make up more than 70% of all motorized trips in Shenzhen.
Ridership
Since the opening of the first phase in 2004, there has been a steady growth in passenger traffic. In 2009 and 2010, passenger traffic soared with major openings of new phase 2 lines, with a three-fold increase in passenger traffic in 2010. On July 12, 2019, it set a new record for its peak ridership at 6.63 million.
July is the busiest month of the year for the Shenzhen Metro, accounting for 9.3% of annual passenger traffic, while January is the least busy month, accounting for only 6.7%. This is caused by Shenzhen's large migrant worker population.
Fares and tickets
Metro rides are priced according to distance travelled, and fares vary from 2 RMB to 14 RMB. Since December 2010 fares are based on a usage fee (2 RMB) + a distance fee. The distance fee is 1 RMB for each from to ; after that 1 RMB for each from to and finally 1 RMB for every over distance. For passengers who wish to ride on business coach in line 11, they have to pay 3 times the amount of price that calculated by the regulations above.
Children under the height of or aged below 6 may ride for free when accompanied by an adult. The metro also offers free rides to senior citizens over the age of 65, the physically disabled and military personnel. Tickets for children between and , or aged between 6 and 14 years, or middle school students, are half priced.
Metro fares can be paid for with single-ride tokens, multiple-ride Shenzhen Tong cards or 1- day passes.
Tokens
When using cash, a RFID token (NXP Mifare Classic) is purchased and used for a single, non-returnable journey. There are two different types of tokens, with green tokens for Standard Class, and yellow tokens used for Business Class which is only available on Line 11. All ticket vending machines offer both English and Chinese interface. The purchaser touches a station name to calculate the fare. After payment, a green token is dispensed, which must be scanned at the entrance station and deposited at the exit station. A penalty applies should a token be lost. Purchasers of green tokens cannot ride Business Class on Line 11 directly. Instead, they must get off at any transfer stations with Line 11 and purchase a separate yellow token.
Note that as of 2015, many machines accept only 5 or 10 RMB notes. The token(s) are only valid at the station where issued. Passengers are unable to buy an extra token for return journey prior to departure. Baggage X-Ray machines are located at each station, and may be staffed during peak hours.
Shenzhen Tong cards
Shenzhen Tong is a pre-paid currency card similar to Oyster card system in London and the Octopus card system used in Hong Kong. The multiple fare card stores credit purchased at stations. The card can be used by waving it in front of the card reader located at all entrances and exits to the subway system. Riders who pay for metro fare with a card receive a 5% discount. Since March 1, 2008, riders who pay for a bus fare with a card and then a subway fare within 90 minutes receive an additional 0.4 RMB discount on the subway fare. Card users pay a distance based fare.
Since June 30, 2011, cards containing both a Shenzhen Tong and Hong Kong Octopus chip have been available in both Shenzhen and Hong Kong. There are plans to further integrate the two systems, and for a new card which will be accepted all over Guangdong province and China's two SARs.
Unlike Hong Kong Octopus Cards, Shenzhen Tong cards cannot be sold back to the stations or have faults dealt with by SZMC. Instead, the customer must go to the offices of Shenzhen Tong. Students studying in Shenzhen can use the Shenzhen Tong to receive a 50% discount.
Note that discounts are not applicable for people who ride Business Class carriages on Line 11.
Metro cards can also be used on Shenzhen's public bus system.
Metro 1-day passes
Metro 1-day pass is a smart card that allowed the card holder have unlimited access of the metro system in 24 continuous hours. Passengers can purchase a 1-day pass for RMB 25 in the service center in any metro station. The pass will be activated and the passenger will have 24 continuous hour for unlimited access after the first entrance. When the pass expired, the pass is no longer available for entering a station but able to exiting a station and finish a journey in 27.5 hours. The 1-day passes are not applicable for Business Class carriages on Line 11.
Station facilities, amenities and services
Some stations have toilets (free of charge), and public telephones. SZMC also operates luggage storage facilities in the concourse above Luohu Station. Mobile phone service is available throughout the system provided by China Mobile, China Telecom, and China Unicom.
Like the Hong Kong MTR, Guangzhou, and Foshan metros, station announcements are in Mandarin, Cantonese and English. Some announcements, such as train arrival, are in Mandarin and English only. Cantonese, an important local language, is chosen for the local Cantonese population as well as Cantonese speakers in Guangdong, Hong Kong and Macau.
The stations of Line 6 and Line 10 are the first metro stations in China to have 5G coverage.
Equipment
Rolling stock
Line 1
22 Bombardier Transportation Movia 456 6-car sets (101–122, Some trains now use CRRC Times Electric traction units after have a maintenance in Changchun)
4 Changchun Railway Vehicles Type A 6-car sets, traction units by Bombardier Transportation (123–126)
26 Zhuzhou Electric Locomotive Works Type A 6-car sets, traction units by Siemens (127–152, Some trains now use CRRC Times Electric traction units after have a maintenance in Zhuzhou)
33 Zhuzhou Electric Locomotive Works Type A 6-car sets, traction units by CSR Times Electric (153–185)
Line 2
35 Changchun Railway Vehicles Type A 6-car sets, traction units by Bombardier Transportation (201–235)
17 Changchun Railway Vehicles Type A 6-car sets, traction units by Bombardier Transportation (236–239、241–252), by CRRC Qingdao Sifang (240).
5 Zhuzhou Electric Locomotive Works Type A 6-car sets, traction units by CRRC Times Electric (253–257)
Line 3
43 Changchun Railway Vehicles Type B 6-car sets, traction units by Hyundai Rotem (301–343, Some trains now use CRRC Times Electric traction units after have a maintenance in Changchun)
33 Nanjing Puzhen Rolling Stock Works Type B 6-car sets, traction units by Hyundai Rotem (344–376)
22 Nanjing Puzhen Rolling Stock Works Type B 6-car sets, traction units by Jiangsu Kingway Rail (377-398)
Line 4
28 Nanjing Puzhen Rolling Stock Works Type A 6-car sets, traction units by ABB (401–428)
24 Nanjing Puzhen Rolling Stock Works Type A 6-car sets, traction units by Nanjing Huashi (429–452)
Line 5
22 Zhuzhou Electric Locomotive Works Type A 6-car sets, traction units by Siemens (501–522)
8 Zhuzhou Electric Locomotive Works Type A 6-car sets, traction units by CSR Times Electric (523–530)
21 Changchun Railway Vehicles Type A 6-car sets, traction units by Bombardier Transportation (531–551)
6 Zhuzhou Electric Locomotive Works Type A 6-car sets, traction units by CRRC Times Electric (552–557)
Line 6
51 Nanjing Puzhen Rolling Stock Works Type A 6-car sets, traction units by Bombardier Transportation (601–651)
Line 6 Branch
9 Nanjing Puzhen Rolling Stock Works Type B 6-car sets, traction units by INVT (6Z01–6Z09)
Line 7
41 Changchun Railway Vehicles Type A 6-car sets, traction units by Bombardier Transportation (701–741)
Line 8
24 Zhuzhou Electric Locomotive Works Type A 6-car sets, traction units by CRRC Times Electric (258–281[801–824], 273-281[816-824] traction units sounds are much higher)
Line 9
29 Changchun Railway Vehicles Type A 6-car sets, traction units by CRRC Times Electric (901–929).
22 Changchun Railway Vehicles Type A 6-car sets, traction units by INVT (930–951).
Line 10
35 Changchun Railway Vehicles Type A 8-car sets, traction units by CRRC Times Electric (1001–1035)
Line 11
33 Zhuzhou Electric Locomotive Works Type A 8-car sets, traction units by CRRC Times Electric (1101–1133)
40 Changchun Railway Vehicles Type A 8-car sets, traction units by CRRC Times Electric (1134–1173)
Line 12
56 Nanjing Puzhen Rolling Stock Works Type A 6-car sets, traction units by Bombardier Transportation (1201–1256)
Line 13
19 CRRC Qingdao Sifang Type A 8-car sets, traction units by Siemens (1301-1319)
Line 14
44 Changchun Railway Vehicles Type A 8-car sets, traction units by CRRC Times Electric (1401–1444)
Line 16
32 Zhuzhou Electric Locomotive Works Type A 6-car sets, traction units by Jiangsu Kingway Rail (1601–1632)
Line 20
9 Changchun Railway Vehicles Type A 8-car sets, traction units by CRRC Times Electric (2001–2009)
Signalling system
On Line 1 and Line 4, Siemens supplied 7 (Phase 1) and 6 (Phase 2) LZB 700 M continuous automatic control systems; 7 (Phase 1) and 6 (Phase 2) electronic Sicas ESTT interlockings; the Vicos OC 501 operations control system with 2 operations control centers, fall-back level with Vicos OC 101 and RTU (FEP), 230 (Phase 1) and 240 (Phase 2) FTG S track vacancy detection units.
Line 2 and Line 5 use Casco CBTC system with 2.4 GHz frequencies, and so the system has suffered frequent problems with interference from consumer Wi-Fi equipment. By the end of November 2012, CASCO solved the problem on Lines 2 and 5 by switching to their standard solution with frequency diversity on 2 different channels.
Accidents and incidents
April 4, 2011 – One worker was killed and four others injured on April 4 when a manually controlled chain hoist broke loose in a Line 5 tunnel in Longgang district. A preliminary investigation by district safety authorities found mechanical failure was to blame.
September 9, 2013 – Three passengers abandoned in Line 1 tunnel after train door opens.
June 25, 2015 – Worker killed during tunnel collapse in Line 7 construction.
April 19, 2017 – Scaffolding for a metro station collapsed during the construction of the Line 8 on Yantian Rd, killing a worker and injuring three.
May 11, 2017 – During the construction of the extension of Line 3 heavy rains caused a partial cave in at an excavation pit for a station on the southern extension of Line 3, killing 2 workers and injuring another.
July 5–7, 2018 – Over a span of three days, at least seven incidents occurred, where power cables were accidentally cut at various construction sites of Shenzhen Metro, causing blackouts in large areas.
July 10, 2018 – During the construction of Line 10, workers accidentally dug up the pipes of Shenzhen Buji Water Supply Co., Ltd., disrupting water distribution system. The Shenzhen Economic and Information Commission warned and penalized the contractor responsible.
Network map
Connections
Pingshan Center station on Line 14 connects to Line 1 of the Pingshan Skyshuttle. This 8.5 km monorail line opened on December 28, 2022, with 11 stations, all in Pingshan district. Two of the stations will eventually connect to Line 16 stations.
Branding
The mascot of Shenzhen Metro named "Tiebao" (Simplified Chinese: 铁宝; Traditional Chinese: 鐵寳) was unveiled on December 27, 2023. It is a cute anthropomorphic of a Shenzhen Metro train's front.
| Technology | China | null |
1374448 | https://en.wikipedia.org/wiki/Degree%20%28graph%20theory%29 | Degree (graph theory) | In graph theory, the degree (or valency) of a vertex of a graph is the number of edges that are incident to the vertex; in a multigraph, a loop contributes 2 to a vertex's degree, for the two ends of the edge. The degree of a vertex is denoted or . The maximum degree of a graph is denoted by , and is the maximum of 's vertices' degrees. The minimum degree of a graph is denoted by , and is the minimum of 's vertices' degrees. In the multigraph shown on the right, the maximum degree is 5 and the minimum degree is 0.
In a regular graph, every vertex has the same degree, and so we can speak of the degree of the graph. A complete graph (denoted , where is the number of vertices in the graph) is a special kind of regular graph where all vertices have the maximum possible degree, .
In a signed graph, the number of positive edges connected to the vertex is called positive deg and the number of connected negative edges is entitled negative deg.
Handshaking lemma
The degree sum formula states that, given a graph ,
.
The formula implies that in any undirected graph, the number of vertices with odd degree is even. This statement (as well as the degree sum formula) is known as the handshaking lemma. The latter name comes from a popular mathematical problem, which is to prove that in any group of people, the number of people who have shaken hands with an odd number of other people from the group is even.
Degree sequence
The degree sequence of an undirected graph is the non-increasing sequence of its vertex degrees; for the above graph it is (5, 3, 3, 2, 2, 1, 0). The degree sequence is a graph invariant, so isomorphic graphs have the same degree sequence. However, the degree sequence does not, in general, uniquely identify a graph; in some cases, non-isomorphic graphs have the same degree sequence.
The degree sequence problem is the problem of finding some or all graphs with the degree sequence being a given non-increasing sequence of positive integers. (Trailing zeroes may be ignored since they are trivially realized by adding an appropriate number of isolated vertices to the graph.) A sequence which is the degree sequence of some simple graph, i.e. for which the degree sequence problem has a solution, is called a graphic or graphical sequence. As a consequence of the degree sum formula, any sequence with an odd sum, such as (3, 3, 1), cannot be realized as the degree sequence of a graph. The inverse is also true: if a sequence has an even sum, it is the degree sequence of a multigraph. The construction of such a graph is straightforward: connect vertices with odd degrees in pairs (forming a matching), and fill out the remaining even degree counts by self-loops.
The question of whether a given degree sequence can be realized by a simple graph is more challenging. This problem is also called graph realization problem and can be solved by either the Erdős–Gallai theorem or the Havel–Hakimi algorithm.
The problem of finding or estimating the number of graphs with a given degree sequence is a problem from the field of graph enumeration.
More generally, the degree sequence of a hypergraph is the non-increasing sequence of its vertex degrees. A sequence is -graphic if it is the degree sequence of some simple -uniform hypergraph. In particular, a -graphic sequence is graphic. Deciding if a given sequence is -graphic is doable in polynomial time for via the Erdős–Gallai theorem but is NP-complete for all .
Special values
A vertex with degree 0 is called an isolated vertex.
A vertex with degree 1 is called a leaf vertex or end vertex or a pendant vertex, and the edge incident with that vertex is called a pendant edge. In the graph on the right, {3,5} is a pendant edge. This terminology is common in the study of trees in graph theory and especially trees as data structures.
A vertex with degree n − 1 in a graph on n vertices is called a dominating vertex.
Global properties
If each vertex of the graph has the same degree k, the graph is called a k-regular graph and the graph itself is said to have degree k. Similarly, a bipartite graph in which every two vertices on the same side of the bipartition as each other have the same degree is called a biregular graph.
An undirected, connected graph has an Eulerian path if and only if it has either 0 or 2 vertices of odd degree. If it has 0 vertices of odd degree, the Eulerian path is an Eulerian circuit.
A directed graph is a directed pseudoforest if and only if every vertex has outdegree at most 1. A functional graph is a special case of a pseudoforest in which every vertex has outdegree exactly 1.
By Brooks' theorem, any graph G other than a clique or an odd cycle has chromatic number at most Δ(G), and by Vizing's theorem any graph has chromatic index at most Δ(G) + 1.
A k-degenerate graph is a graph in which each subgraph has a vertex of degree at most k.
| Mathematics | Graph theory | null |
34360942 | https://en.wikipedia.org/wiki/Graecopithecus | Graecopithecus | Graecopithecus is an extinct genus of hominid that lived in southeast Europe during the late Miocene around 7.2 million years ago. Originally identified by a single lower jawbone bearing teeth found in Pyrgos Vasilissis, Athens, Greece, in 1944, other teeth were discovered from Azmaka quarry in Bulgaria in 2012. With only little and badly preserved materials to reveal its nature, it is considered as "the most poorly known European Miocene hominoids." The creature was popularly nicknamed 'El Graeco' (word play on the Greek-Spanish painter El Greco) by scientists.
In 2017, an international team of palaeontologists led by Madelaine Böhme of the Eberhard-Karls-University Tübingen, Germany, published a controversial analysis of the teeth and age of the specimens, and came to the conclusion that it could be the oldest hominin, meaning that it could be the oldest direct ancestors of humans after splitting from that of the chimpanzees. Their simultaneous study also claimed that contrary to the generally accepted evidence of the African origin of the hominin lineage, the ancestors of humans originated from the main ape ancestry in the Mediterranean region (before migrating into Africa where they evolved into the ancestors of Homo species). They named the origin of human theory as the "North Side Story."
These claims have been disputed by other scientists. Rick Potts and Bernard Wood argued that the evidence is too flimsy to even say it is a hominin. Tim D. White commented that the claim was only to support a biased argument that Africa is not the birthplace of humans; while Sergio Almécija stated that single characters such as teeth cannot tell the claimed evolutionary details.
Discovery
The original Graecopithecus specimen was a single lower jawbone (mandible) found from a site called Pyrgos Vassilissis, northwest of Athens, in southern Greece in 1944, "reportedly unearthed as the occupying German forces were building a wartime bunker". The jawbone was almost complete with teeth when it was sent to Berlin for analysis, but was damaged by bombings during the final phases of World War II. Only the second molar and fourth premolar remain intact, while fragments of other teeth are still embedded. The original finder, German paleontologist Bruno von Freyberg initially believed that it belonged to an extinct Old World monkey Mesopithecus, as he reported in 1951. However, Gustav Heinrich Ralph von Koenigswald realised that it was the tooth of an ape family and erected the scientific name Graecopithecus freybergi in 1972, after the discoverer.
Another tooth remain was discovered from Azmaka quarry in Bulgaria in 2012.
Description
The mandible of Graecopithecus with a third molar that is very worn, the root of a second molar, and a fragment of a premolar, is dated from the late Miocene around 7.2 million years old. Excavation of the site is not possible (as of 1986) due to the owner having built a swimming pool on the location.
The thick enamel and large molars are the features that convinced von Koenigswald that the specimen belonged to a hominid species. X-ray microtomography and 3-dimensional reconstruction in 2017 revealed that it belonged to an adult individual and possibly a male. The partial fusion of the fourth premolar (P4) roots is an additional evidence that it is of a hominid, and the thick enamel resembles those of the human lineage (hominins).
Classification
G. freybergi is considered to be possibly the same taxon as Ouranopithecus macedoniensis, another extinct hominid described in 1977 from northern Greece. Due to paucity of specimens and poor quality of the fossils, it remains the least well-known extinct hominid found within Europe. In 1984, British palaeontologists Peter Andrews and Lawrence B. Martin classified Graecopithecus and Ouranopithecus as synonyms (same taxon) and treated them as members of the genus Sivapithecus. This classification persisted for several years until additional Ouranopithecus fossils were discovered including part of the skull in the 1990s that indicated better distinction as different hominids. Based on new evidences, in 1997, Australian palaeontologist David W. Cameron proposed renaming and inclusion of Ouranopithecus into Graecopithecus based on taxonomic priority with Graecopithecus macedoniensis as a new name for O. macedoniensis. However, better O. macedoniensis specimens were found including a new species Ouranopithecus turkae from Turkey that supported separation of the genus. This change was generally adopted.
Re-examination and reinterpretation
In 2017, an international team of palaeontologists led by Madelaine Böhme (Eberhard-Karls-University Tübingen, Germany) published detailed reanalysis and new interpretation in the journal PLOS One. One paper deals with an examination of the detailed morphology of molar teeth of G. freybergi from Greece and Bulgaria, and compared it with that of Ouranopithecus. The study concluded that Graecopithecus was a hominin, sharing ancestry with Homo but not with the chimpanzees (Pan), and distinct from Ouranopithecus, which has more ape-like traits. If this classification is correct, Graecopithecus would be the oldest known representative of the human lineage after the human–chimpanzee split, in 19th-century terminology, the "missing link" between human and non-human primates. The species was found to be some two hundred thousand years older than the oldest known hominid found in Africa (not necessarily ancestral to the human lineage), Sahelanthropus tchadensis. The study concludes:[The] dental root attributes of Graecopithecus suggest hominin affinities, such that its hominin status cannot be excluded. If this status is confirmed by additional fossil evidence, Graecopithecus would be the oldest known hominin and the oldest-known crown hominine, as the evidence for the gorillin status of Chororapithecus is much weaker than the hominin status of Graecopithecus. More fossils are needed but at this point it seems likely that the Eastern Mediterranean needs to be considered as just as likely a place of hominine diversification and hominin origins as tropical Africa.An accompanying paper presents the study of the geological environments of the areas where the fossils were discovered. Until then, the precise date of Graecopithecus has not been resolved and usually inferred from geological data of materials related the fossils and surrounding areas that add to uncertainty in its evolutionary importance and relationship with other hominids. It is often broadly described as 6.6 to 8 million years old. The PLOS One paper resolved that the hominid lived 7.37 to 7.11 million years ago, with the specime from Greece dated to 7.18 Ma and that from Bulgaria to 7.24 Ma. It also indicates that as the species lived in Europe, it suggest "that major splits in the hominid family occurred outside Africa."
It has also been proposed the Graecopithecus may not be a direct ancestor of the human lineage, but instead may have evolved its hominin-like traits independently. The emergence of Homo itself is dated to close to 4 million years later than Graecopithecus, so that the appearance of Graecopithecus in Europe does not preclude the development of Homo proper in East Africa (as suggested by Homo habilis being found in Tanzania); however, the popular press reporting on the 2017 study did cast its result in terms of determining the "birthplace of mankind". Graecopithecus lived in southeast Europe 7.2 million years ago, and if the premise of the study is correct, Graecopithecus, after evolving in Europe, would have migrated to Africa about 7 million years ago where its descendants would eventually evolve into the genus Homo.
Criticism
The 2017 PLOS One papers made two critical conclusions: that Graecopithecus is a hominin suggesting it as the oldest ancestor of humans after splitting from chimpazees, and that as Graecopithecus is a human ancestor, Europe is the birthplace of hominins. This directly challenges the prevailing knowledge that humans originated in East Africa.
David R. Begun of the University of Toronto, Canada, one of the co-authors, was quoted as saying that "[t]his dating allows us to move the human–chimpanzee split into the Mediterranean area." This was set against a quote by an uninvolved anthropologist saying that "[i]t is possible that the human lineage originated in Europe, but very substantial fossil evidence places the origin in Africa [...] I would be hesitant about using a single character from an isolated fossil to set against the evidence from Africa." Since 1994, Begun had adhered to a hypothesis that African hominids (including living apes) descended from Eurasian apes since the older ape fossils are found in Europa and Asia. This is a feasible explanation as it is possible that the African ape ancestors could move to Africa around 9 million years ago from Europe.
However, claiming that Graecopithecus is an evidence of human origin in Europe is illogical since all human ancestral species known so far are strictly found in Africa; as Rick Potts, head of the Smithsonian's Human Origins Program, remarked: "I think the principal claim of the main paper goes well beyond the evidence in hand... A hominin or even a hominine (modern African ape) ancestor located in a fairly isolated place in southern Europe doesn’t make much sense geographically as the ancestor of modern African apes, or particular the oldest ancestor of African hominins." David Alba at the Catalan Institute of Palaeontology in Barcelona was the first to point out that "It is not surprising at all that Begun is now arguing that hominins as well originated in Europe." Julien Benoit of the University of the Witwatersrand, Johannesburg, South Africa, also commented: "Any study that counters this consensus (Out of Africa theory) would have to provide very strong evidence and perfect methodology to support its claim. In my opinion, this article doesn't meet those criteria."
Other scientists have also expressed skepticism of Begun's classification. Bernard Wood at George Washington University described the hypothesis as "relatively weak" and Sergio Almécija, also at George Washington University, says it is important to bear in mind that primates seem particularly prone to evolving similar features independently. "Single characters are not reliable to make big evolutionary [claims]." Tim White at the University of California, Berkeley, asserted that the study was merely an attempt "to resurrect Begun’s tired argument with a long-known crappy fossil, newly scanned."
Reassessment
In late 2017, Julien Benoit and Francis J. Thackeray re-analysed the claims of the PLOS One papers and found key issues in the major conclusions:
The partial fusion of the fourth premolar (P4) roots does not define Graecopithecus as a hominin since the feature is common in hominids, even in the chimpanzees.
Thick enamel and relatively large molar teeth are not exclusive to hominins as they are also present in other Miocene apes and gorillas.
The claim that Graecopithecus is the ancestral ape of human lineage and that humans originate in Europe is not justified. Even if Graecopithecus is the basal (root ancestor) ape, all other human ancestral species starting from Sahelanthropus were in Africa, thus, still making Africa the birthplace of humans.
The study concludes:[We] recognise a small signal for placing Graecopithecus at the root of the Hominini clade. This means that the phylogenetic relationship between Graecopithecus and Hominini is as yet not confirmed. Our analysis supports the view that Graecopithecus is potentially an important taxon for the origin of Hominini, but this is not certain and deserves further investigation and more material.
Response
In 2018, Fuss, Spassov, Böhme, and Begun published a response to Benoit and Thackeray, claiming that their original publication had been misrepresented and misconstrued. They explained that the conclusion of the 2017 paper had not been that Graecopithecus was certainly a hominin, but that its status as a hominin could not be ruled out, and that more research and evidence would be needed to make a conclusion—a conclusion that Benoit and Thackeray make in their own paper. They argued against Benoit and Thackeray write that they did not judge canine root derivation of Graecopithecus and Salehanthropus against each other, stating that the differences between them were within the range of sexual variation. Additionally, when Benoit and Thackeray claim that the characteristics mentioned in the 2017 paper are not unique to Hominini, they do not mention that the 2017 paper discusses canine root size and premolar root complexity reduction, which could be indications of Hominini.
| Biology and health sciences | Australopithecines | Biology |
27533725 | https://en.wikipedia.org/wiki/Intergalactic%20star | Intergalactic star | An intergalactic star, also known as an intracluster star or a rogue star, is a star not gravitationally bound to any galaxy. Although a source of much discussion in the scientific community during the late 1990s, intergalactic stars are now generally thought to have originated in galaxies, like other stars, before being expelled as the result of either galaxies colliding or of a multiple-star system traveling too close to a supermassive black hole, which are found at the center of many galaxies.
Collectively, intergalactic stars are referred to as the intracluster stellar population, or IC population for short, in the scientific literature.
Discovery
The hypothesis that stars exist only in galaxies was disproven in 1997 with the discovery of intergalactic stars. The first to be discovered were in the Virgo cluster of galaxies, where some one trillion are now surmised to exist.
Formation
The way these stars arise is still a mystery, but several scientifically credible hypotheses have been suggested and published by astrophysicists.
The most common hypothesis is that the collision of two or more galaxies can toss some stars out into the vast empty regions of intergalactic space. Although stars normally reside within galaxies, they can be expelled by gravitational forces when galaxies collide. It is commonly believed that intergalactic stars may primarily have originated from extremely small galaxies, since it is easier for stars to escape a smaller galaxy's gravitational pull, than that of a large galaxy. However, when large galaxies collide, some of the gravitational disturbances might also expel stars. In 2015, a study of supernovae in intergalactic space suggested that the progenitor stars had been expelled from their host galaxies during a galactic collision between two giant ellipticals, as their supermassive black hole centres merged.
Another hypothesis, that is not mutually exclusive to the galactic collisions hypothesis, is that intergalactic stars were ejected from their galaxy of origin by a close encounter with the supermassive black hole in the galaxy center, should there be one. In such a scenario, it is likely that the intergalactic star(s) was originally part of a multiple star system where the other stars were pulled into the supermassive black hole and the soon-to-be intergalactic star was accelerated and ejected away at very high speeds. Such an event could theoretically accelerate a star to such high speeds that it becomes a hypervelocity star, thereby escaping the gravitational well of the entire galaxy. In this respect, model calculations (from 1988) predict the supermassive black hole in the center of our Milky Way galaxy to expel one star every 100,000 years on average.
Observation history
In 1997, the Hubble Space Telescope discovered a large number of intergalactic stars in the Virgo cluster of galaxies. Later in the 1990s, scientists discovered another group of intergalactic stars in the Fornax cluster of galaxies.
In 2005, at the Smithsonian Center for Astrophysics, Warren Brown and his team attempted to measure the speeds of hypervelocity stars by using the Doppler Technique, by which light is observed for the similar changes that occur in sound when an object is moving away or toward something. But the speeds found are only estimated minimums, as in reality their speeds may be larger than the speeds found by the researchers. "One of the newfound exiles is moving in the direction of the constellation Ursa Major at about 1.25 million mph with respect to the galaxy. It is 240,000 light-years away. The other is headed toward the constellation Cancer, outbound at 1.43 million miles per hour and 180,000 light-years away."
In the late 2000s, a diffuse glow from the intergalactic medium, but of unknown origin, was discovered. In 2012, it was suggested and shown that it might originate from intergalactic stars. Subsequent observations and studies have elaborated on the issue and described the diffuse extragalactic background radiation in more detail.
Some Vanderbilt astronomers report that they have identified more than 675 stars at the edge of the Milky Way, between the Andromeda Galaxy and the Milky Way. They argue that these stars are hypervelocity (intergalactic) stars that were ejected from the Milky Way's Galactic Center. These stars are red giants with a high metallicity (a measure of the proportion of chemical elements other than hydrogen and helium within a star) indicating an inner galactic origin, since stars outside the disks of galaxies tend to have low metallicity and are older.
Some recently discovered supernovae have been confirmed to have exploded hundreds of thousands of light-years from the nearest star or galaxy. Most intergalactic star candidates found in the neighborhood of the Milky Way seem not to have an origin in the Galactic Center but in the Milky Way disk or elsewhere.
Mass
In 2005, the Spitzer Space Telescope revealed a hitherto unknown infrared component in the background from the cosmos. Since then, several other anisotropies at other wavelengths including blue and x-ray have been detected with other space telescopes and they are now collectively described as the diffuse extragalactic background radiation. Several explanations have been discussed by scientists, but in 2012, it was suggested and shown how for the first time this diffuse radiation might originate from intergalactic stars. If that is the case, they might collectively comprise as much mass as that found in the galaxies. A population of such magnitude was at one point thought to explain the photon underproduction crisis, and may explain a significant part of the dark matter problem.
Known locations
The first intergalactic stars were discovered in the Virgo cluster of galaxies. These stars are notable for their isolation, residing approximately 300,000 light-years away from the nearest galaxy. Despite the difficulty in determining their exact mass, it is estimated that intergalactic stars constitute around 10 percent of the mass of the Virgo cluster, potentially outweighing any of its 2,500 galaxies
In 2012, astronomers identified approximately 675 rogue stars at the edge of the Milky Way, towards the Andromeda Galaxy. These stars were likely ejected from the Milky Way's core by interactions with the central supermassive black hole. The study led by Kelly Holley-Bockelmann and Lauren Palladino from Vanderbilt University highlighted the unusual red coloration and high velocities of these stars, indicating their dramatic journey from the galactic center.
| Physical sciences | Stellar astronomy | Astronomy |
537264 | https://en.wikipedia.org/wiki/Fractional%20freezing | Fractional freezing | Fractional freezing is a process used in process engineering and chemistry to separate substances with different melting points. It can be done by partial melting of a solid, for example in zone refining of silicon or metals, or by partial crystallization of a liquid, as in freeze distillation, also called normal freezing or progressive freezing. The initial sample is thus fractionated (separated into fractions).
Partial crystallization can also be achieved by adding a dilute solvent to the mixture, and cooling and concentrating the mixture by evaporating the solvent, a process called solution crystallization. Fractional freezing is generally used to produce ultra-pure solids, or to concentrate heat-sensitive liquids.
Freeze distillation
Freeze distillation is actually a condensation, not distillation per se. It removes frozen solid, separating a dissolved material from the liquid left behind. This is analogous to true distillation, where the evaporated and re-condensed fraction is enriched.
Ethanol and liquid water are completely miscible, but ethanol is practically insoluble in water ice. That means almost pure water ice can be precipitated from a lean ethanol-water mixture by cooling it sufficiently. The precipitation of water ice from the mixture enriches ethanol in the remaining liquid phase. The two phases can then be separated by filtration or decanting. The temperature at which water ice starts to precipitate depends on the ethanol concentration. Consequently, at a given temperature and ethanol concentration, the freezing process will reach an equilibrium at a specific ratio of water ice and enriched ethanol solution with a specific ethanol concentration. The temperatures and mixing ratios of these phase equilibria can be read from the phase diagram of ethanol and water. The maximum enrichment of ethanol in the liquid phase is reached at the eutectic point of ethanol and water, approximately 92.4 weight-% ethanol at -123 °C.
The best-known freeze-distilled beverages are applejack and ice beer. Ice wine is the result of a similar process, but in this case, the freezing happens before the fermentation, and thus it is sugar, not alcohol, that gets concentrated.
Purification of solids
When a pure solid is desired, two possible situations can occur. If the contaminant is soluble in the desired solid, a multiple stage fractional freezing is required, analogous to multistage distillation. If, however, a eutectic system forms (analogous to an azeotrope in distillation), a very pure solid can be recovered, as long as the liquid is not at its eutectic composition (in which case a mixed solid forms, which can be hard to separate) or above its eutectic composition (in which case the undesired solid forms).
Concentration of liquids
When the requirement is to concentrate a liquid phase, fractional freezing can be useful due to its simplicity. Fractional freezing is also used in the production of fruit juice concentrates and other heat-sensitive liquids, as it does not involve heating the liquid (as happens during evaporation).
Desalination
Fractional freezing can be used to desalinate sea water. In a process that naturally occurs with sea ice, frozen salt water, when partially melted, leaves behind ice that is of a much lower salt content. Because sodium chloride lowers the melting point of water, the salt in sea water tends to be forced out of pure water while freezing, called brine rejection. Large lakes of higher salinity water, called polynyas, form in the middle of floes, and the water eventually sinks to the bottom. Likewise, the frozen water with the highest concentration of salt melts first. Either method decreases the salinity of the remaining frozen water, and after multiple runs the results can be drinkable.
Alcoholic beverages
Fractional freezing can be used as a simple method to increase the alcohol concentration in fermented alcoholic beverages, a process sometimes called freeze distillation. Examples are applejack, made from hard cider, and ice beer. In practice, while not able to produce an alcohol concentration comparable to distillation, this technique can achieve some concentration with far less effort than any practical distillation apparatus would require. The danger of fractional freezing of alcoholic beverages is that it does not remove impurities, unlike (heat) distillation. Thus, fractional freezing may increase the ratio of impurities to the total volume of the beverage (though not necessarily the ratio of impurities to the amount of ethanol). This concentration may cause side effects to the drinker, leading to intense hangovers and a condition known as "apple palsy" (although this term has also simply been used to refer to intoxication, especially from applejack.)
Alternative fuels
Fractional freezing is commonly used as a simple method to reduce the gel point of biodiesel and other alternative diesel fuels, whereby esters of higher gel point are removed from esters of lower gel point through cold filtering, or other methods to reduce the subsequent alternative fuel gel point of the fuel blend. This process employs fuel stratification whereby components in the fuel blend develop a higher specific gravity as they approach their respective gel points and thus sink to the bottom of the container, where they can be removed.
| Physical sciences | Phase separations | Chemistry |
537442 | https://en.wikipedia.org/wiki/HDMI | HDMI | High-Definition Multimedia Interface (HDMI) is a proprietary audio/video interface for transmitting uncompressed video data and compressed or uncompressed digital audio data from a source device, such as a display controller, to a computer monitor, video projector, digital television, or digital audio device. HDMI is a digital replacement for analog video standards.
HDMI implements the ANSI/CTA-861 standard, which defines video formats and waveforms, transport of compressed and uncompressed LPCM audio, auxiliary data, and implementations of the VESA EDID. CEA-861 signals carried by HDMI are electrically compatible with the CEA-861 signals used by the Digital Visual Interface (DVI). No signal conversion is necessary, nor is there a loss of video quality when a DVI-to-HDMI adapter is used. The Consumer Electronics Control (CEC) capability allows HDMI devices to control each other when necessary and allows the user to operate multiple devices with one handheld remote control device.
Several versions of HDMI have been developed and deployed since the initial release of the technology, occasionally introducing new connectors with smaller form factors, but all versions still use the same basic pinout and are compatible with all connector types and cables. Other than improved audio and video capacity, performance, resolution and color spaces, newer versions have optional advanced features such as 3D, Ethernet data connection, and CEC extensions.
Production of consumer HDMI products started in late 2003. In Europe, either DVI-HDCP or HDMI is included in the HD ready in-store labeling specification for TV sets for HDTV, formulated by EICTA with SES Astra in 2005. HDMI began to appear on consumer HDTVs in 2004 and camcorders and digital still cameras in 2006. , nearly 10 billion HDMI devices have been sold.
History
The HDMI founders were Hitachi, Panasonic, Philips, Silicon Image, Sony, Thomson, and Toshiba. Digital Content Protection, LLC provides HDCP (which was developed by Intel) for HDMI. HDMI has the support of motion picture producers Fox, Universal, Warner Bros. and Disney, along with system operators DirecTV, EchoStar (Dish Network) and CableLabs.
The HDMI founders began development on HDMI 1.0 on April 16, 2002, with the goal of creating an AV connector that was backward-compatible with DVI. At the time, DVI-HDCP (DVI with HDCP) and DVI-HDTV (DVI-HDCP using the CEA-861-B video standard) were being used on HDTVs. HDMI 1.0 was designed to improve on DVI-HDTV by using a smaller connector and adding audio capability and enhanced capability and consumer electronics control functions.
The first Authorized Testing Center (ATC), which tests HDMI products, was opened by Silicon Image on June 23, 2003, in California, United States. The first ATC in Japan was opened by Panasonic on May 1, 2004, in Osaka. The first ATC in Europe was opened by Philips on May 25, 2005, in Caen, France. The first ATC in China was opened by Silicon Image on November 21, 2005, in Shenzhen. The first ATC in India was opened by Philips on June 12, 2008, in Bangalore. The HDMI website contains a list of all the ATCs.
According to In-Stat, the number of HDMI devices sold was 5 million in 2004, 17.4 million in 2005, 63 million in 2006, and 143 million in 2007. HDMI has become the de facto standard for HDTVs, and according to In-Stat, around 90% of digital televisions in 2007 included HDMI. In-Stat has estimated that 229 million HDMI devices were sold in 2008. On April 8, 2008, there were over 850 consumer electronics and PC companies that had adopted the HDMI specification (HDMI adopters). On January 7, 2009, HDMI Licensing, LLC announced that HDMI had reached an installed base of over 600 million HDMI devices. In-Stat estimated that 394 million HDMI devices would sell in 2009 and that all digital televisions by the end of 2009 would have at least one HDMI input.
On January 28, 2008, In-Stat reported that shipments of HDMI were expected to exceed those of DVI in 2008, driven primarily by the consumer electronics market.
In 2008, PC Magazine awarded a Technical Excellence Award in the Home Theater category for an "innovation that has changed the world" to the CEC portion of the HDMI specification. Ten companies were given a Technology and Engineering Emmy Award for their development of HDMI by the National Academy of Television Arts and Sciences on January 7, 2009.
On October 25, 2011, the HDMI Forum was established by the HDMI founders to create an open organization so that interested companies can participate in the development of the HDMI specification. All members of the HDMI Forum have equal voting rights, may participate in the Technical Working Group, and if elected can be on the Board of Directors. There is no limit to the number of companies allowed in the HDMI Forum though companies must pay an annual fee of US$15,000 with an additional annual fee of $5,000 for those companies that serve on the Board of Directors. The Board of Directors is made up of 11 companies who are elected every two years by a general vote of HDMI Forum members. All future development of the HDMI specification take place in the HDMI Forum and are built upon the HDMI 1.4b specification. Also on the same day HDMI Licensing, LLC announced that there were over 1,100 HDMI adopters and that over 2 billion HDMI-enabled products had shipped since the launch of the HDMI standard. From October 25, 2011, all development of the HDMI specification became the responsibility of the newly created HDMI Forum.
On January 8, 2013, HDMI Licensing, LLC announced that there were over 1,300 HDMI adopters and that over 3 billion HDMI devices had shipped since the launch of the HDMI standard. The day also marked the 10th anniversary of the release of the first HDMI specification.
, nearly 10 billion HDMI devices had been sold.
Specifications
The HDMI specification defines the protocols, signals, electrical interfaces and mechanical requirements of the standard. The maximum pixel clock rate for HDMI 1.0 is 165 MHz, which is sufficient to allow 1080p and WUXGA (1920×1200) at 60Hz. HDMI 1.3 increases that to 340 MHz, which allows for higher resolution (such as WQXGA, 2560×1600) across a single digital link. An HDMI connection can either be single-link (type A/C/D) or dual-link (type B) and can have a video pixel rate of 25 MHz to 340 MHz (for a single-link connection) or 25 MHz to 680 MHz (for a dual-link connection). Video formats with pixel rates below 25 MHz (like 480i at 13.5 MHz) are transmitted over TMDS links using a pixel-repetition scheme.
Audio/video
HDMI uses the Consumer Technology Association/Electronic Industries Alliance 861 standards. HDMI 1.0 to HDMI 1.2a uses the EIA/CEA-861-B video standard, HDMI 1.3 uses the CEA-861-D video standard, and HDMI 1.4 uses the CEA-861-E video standard. The CEA-861-E document defines "video formats and waveforms; colorimetry and quantization; transport of compressed and uncompressed LPCM audio; carriage of auxiliary data; and implementations of the Video Electronics Standards Association (VESA) Enhanced Extended Display Identification Data Standard (E-EDID)". On July 15, 2013, the CEA announced the publication of CEA-861-F, a standard that can be used by video interfaces such as DVI, HDMI, and LVDS. CEA-861-F adds the ability to transmit several Ultra HD video formats and additional color spaces.
To ensure baseline compatibility between different HDMI sources and displays (as well as backward compatibility with the electrically compatible DVI standard) all HDMI devices must implement the sRGB color space at 8 bits per component. Ability to use the color space and higher color depths ("deep color") is optional. HDMI permits sRGB 4:4:4 chroma subsampling (8–16 bits per component), xvYCC 4:4:4 chroma subsampling (8–16 bits per component), 4:4:4 chroma subsampling (8–16 bits per component), or 4:2:2 chroma subsampling (8–12 bits per component). The color spaces that can be used by HDMI are ITU-R BT.601, ITU-R BT.709-5 and IEC 61966-2-4.
For digital audio, if an HDMI device has audio, it is required to implement the baseline format: stereo (uncompressed) PCM. Other formats are optional, with HDMI allowing up to 8 channels of uncompressed audio at sample sizes of 16 bits, 20 bits, or 24 bits, with sample rates of 32kHz, 44.1kHz, 48kHz, 88.2kHz, 96kHz, 176.4kHz, or 192kHz. HDMI also carries any IEC 61937-compliant compressed audio stream, such as Dolby Digital and DTS, and up to 8 channels of one-bit DSD audio (used on Super Audio CDs) at rates up to four times that of Super Audio CD. With version 1.3, HDMI allows lossless compressed audio streams Dolby TrueHD and DTS-HD Master Audio. As with the video, audio capability is optional. Audio return channel (ARC) is a feature introduced in the HDMI 1.4 standard. "Return" refers to the case where the audio comes from the TV and can be sent "upstream" to the AV receiver using the HDMI cable connected to the AV receiver. An example given on the HDMI website is that a TV that directly receives a terrestrial/satellite broadcast, or has a video source built in, sends the audio "upstream" to the AV receiver.
The HDMI standard was not designed to pass closed caption data (for example, subtitles) to the television for decoding. As such, any closed caption stream must be decoded and included as an image in the video stream(s) prior to transmission over an HDMI cable to appear on the DTV. This limits the caption style (even for digital captions) to only that decoded at the source prior to HDMI transmission. This also prevents closed captions when transmission over HDMI is required for upconversion. For example, a DVD player that sends an upscaled 720p/1080i format via HDMI to an HDTV has no way to pass Closed Captioning data so that the HDTV can decode it, as there is no line 21 VBI in that format.
Communication channels
HDMI has three physically separate communication channels, which are the DDC, TMDS and the optional CEC. HDMI 1.4 added ARC and HEC.
Display Data Channel (DDC)
The Display Data Channel (DDC) is a communication channel based on the I2C bus specification. HDMI specifically requires the device implement the Enhanced Display Data Channel (E-DDC), which is used by the HDMI source device to read the E-EDID data from the HDMI sink device to learn what audio/video formats it can take. HDMI requires that the E-DDC implement I2C standard mode speed () and allows it to optionally implement fast mode speed ().
I2C address 0x74 on the DDC channel is actively used for High-bandwidth Digital Content Protection (HDCP).
Transition-minimized differential signaling (TMDS)
Transition-minimized differential signaling (TMDS) on HDMI interleaves video, audio and auxiliary data using three different packet types, called the video data period, the data island period and the control period. During the video data period, the pixels of an active video line are transmitted. During the data island period (which occurs during the horizontal and vertical blanking intervals), audio and auxiliary data are transmitted within a series of packets. The control period occurs between video and data island periods.
Both HDMI and DVI use TMDS to send 10-bit characters that are encoded using 8b/10b encoding that differs from the original IBM form for the video data period and 2b/10b encoding for the control period. HDMI adds the ability to send audio and auxiliary data using 4b/10b encoding for the data island period. Each data island period is 32 pixels in size and contains a 32-bit packet header, which includes 8 bits of BCH ECC parity data for error correction and describes the contents of the packet. Each packet contains four subpackets, and each subpacket is 64 bits in size, including 8 bits of BCH ECC parity data, allowing for each packet to carry up to 224 bits of audio data. Each data island period can contain up to 18 packets. Seven of the 15 packet types described in the HDMI 1.3a specifications deal with audio data, while the other 8 types deal with auxiliary data. Among these are the general control packet and the gamut metadata packet. The general control packet carries information on AVMUTE (which mutes the audio during changes that may cause audio noise) and color depth (which sends the bit depth of the current video stream and is required for deep color). The gamut metadata packet carries information on the color space being used for the current video stream and is required for xvYCC.
Consumer Electronics Control (CEC)
Consumer Electronics Control (CEC) is an HDMI feature designed to allow the user to command and control up to 15 CEC-enabled devices, that are connected through HDMI, by using only one of their remote controls (for example by controlling a television set, set-top box, and DVD player using only the remote control of the TV). CEC also allows for individual CEC-enabled devices to command and control each other without user intervention.
It is a one-wire bidirectional serial bus that is based on the CENELEC standard AV.link protocol to perform remote control functions. CEC wiring is mandatory, although implementation of CEC in a product is optional. It was defined in HDMI Specification 1.0 and updated in HDMI 1.2, HDMI 1.2a and HDMI 1.3a (which added timer and audio commands to the bus). USB to CEC adapters exist that allow a computer to control CEC-enabled devices.
HDMI Ethernet and Audio Return Channel
Introduced in HDMI 1.4, HDMI Ethernet and Audio Return Channel (HEAC) adds a high-speed bidirectional data communication link (HEC) and the ability to send audio data upstream to the source device (ARC). HEAC utilizes two lines from the connector: the previously unused Reserved pin (called HEAC+) and the Hot Plug Detect pin (called HEAC−). If only ARC transmission is required, a single mode signal using the HEAC+ line can be used, otherwise, HEC is transmitted as a differential signal over the pair of lines, and ARC as a common mode component of the pair.
Audio Return Channel (ARC)
ARC is an audio link meant to replace other cables between the TV and the A/V receiver or speaker system. This direction is used when the TV is the one that generates or receives the video stream instead of the other equipment. A typical case is the running of an app on a smart TV such as Netflix, but reproduction of audio is handled by the other equipment. Without ARC, the audio output from the TV must be routed by another cable, typically TOSLink or RCA, into the speaker system.
HDMI Ethernet Channel (HEC)
HDMI Ethernet Channel technology consolidates video, audio, and data streams into a single HDMI cable, and the HEC feature enables IP-based applications over HDMI and provides a bidirectional Ethernet communication at . The physical layer of the Ethernet implementation uses a hybrid to simultaneously send and receive attenuated 100BASE-TX-type signals through a single twisted pair.
Compatibility with DVI
HDMI is backward compatible with single-link Digital Visual Interface digital video (DVI-D or DVI-I, but not DVI-A or dual-link DVI). No signal conversion is required when an adapter or asymmetric cable is used, so there is no loss of video quality.
From a user's perspective, an HDMI display can be driven by a single-link DVI-D source, since HDMI and DVI-D define an overlapping minimum set of allowed resolutions and frame-buffer formats to ensure a basic level of interoperability. In the reverse case, a DVI-D monitor has the same level of basic interoperability unless content protection with High-bandwidth Digital Content Protection (HDCP) interferes—or the HDMI color encoding is in component color space instead of RGB, which is not possible in DVI. An HDMI source, such as a Blu-ray player, may require an HDCP-compliant display, and refuse to output HDCP-protected content to a non-compliant display. A further complication is that there is a small amount of display equipment, such as some high-end home theater projectors, designed with HDMI inputs but not HDCP-compliant.
Any DVI-to-HDMI adapter can function as an HDMI-to-DVI adapter (and vice versa). Typically, the only limitation is the gender of the adapter's connectors and the gender of the cables and sockets it is used with.
Features specific to HDMI, such as remote-control and audio transport, are not available in devices that use legacy DVI-D signalling. However, many devices output HDMI over a DVI connector (e.g., ATI 3000-series and NVIDIA GTX 200-series video cards), and some multimedia displays may accept HDMI (including audio) over a DVI input. Exact capabilities beyond basic compatibility vary. Adapters are generally bi-directional.
Content protection (HDCP)
High-bandwidth Digital Content Protection (HDCP) is a newer form of digital rights management (DRM). Intel created the original technology to make sure that digital content followed the guidelines set by the Digital Content Protection group.
HDMI can use HDCP to encrypt the signal if required by the source device. Content Scramble System (CSS), Content Protection for Recordable Media (CPRM) and Advanced Access Content System (AACS) require the use of HDCP on HDMI when playing back encrypted DVD Video, DVD Audio, HD DVD and Blu-ray Disc. The HDCP repeater bit controls the authentication and switching/distribution of an HDMI signal. According to HDCP Specification 1.2 (beginning with HDMI CTS 1.3a), any system that implements HDCP must do so in a fully compliant manner. HDCP testing that was previously only a requirement for optional tests such as the "Simplay HD" testing program is now part of the requirements for HDMI compliance. HDCP accommodates up to 127 connected devices with up to 7 levels, using a combination of sources, sinks and repeaters. A simple example of this is several HDMI devices connected to an HDMI AV receiver that is connected to an HDMI display.
Devices called HDCP strippers can remove the HDCP information from the video signal so the video can play on non-HDCP-compliant displays, though a fair use and non-disclosure form must usually be signed with a registering agency before use.
Connectors
There are five HDMI connector types. Type A/B are defined in the HDMI 1.0 specification, type C is defined in the HDMI 1.3 specification, and type D/E are defined in the HDMI 1.4 specification.
Type A; Standard The plug (male) connector outside dimensions are 13.9 mm × 4.45 mm, and the receptacle (female) connector inside dimensions are 14 mm × 4.55 mm. There are 19 pins, with bandwidth to carry all SDTV, EDTV, HDTV, UHD, and 4K modes. It is electrically compatible with single-link DVI-D.
Type B; Dual-link This connector is 21.2 mm × 4.45 mm and has 29 pins, carrying six differential pairs instead of three, for use with very high-resolution displays such as WQUXGA (3840×2400). It is electrically compatible with dual-link DVI-D, but has not yet been used in any products. With the introduction of HDMI 1.3, the maximum bandwidth of single-link HDMI exceeded that of dual-link DVI-D. As of HDMI 1.4, the pixel clock rate crossover frequency from single to dual-link has not been defined.
Type C; Mini This Mini connector is smaller than the type A plug, measuring 10.42 mm × 2.42 mm but has the same 19-pin configuration. It is intended for portable devices. The differences are that all positive signals of the differential pairs are swapped with their corresponding shield, the DDC/CEC Ground is assigned to pin 13 instead of pin 17, the CEC is assigned to pin 14 instead of pin 13, and the reserved pin is 17 instead of pin 14. The type C Mini connector can be connected to a type A connector using a type A-to-type C cable.
Type D; Micro This Micro connector shrinks the connector size to something resembling a micro-USB connector, measuring only 5.83 mm × 2.20 mm For comparison, a micro-USB connector is 6.85 mm × 1.8 mm and a USB type-A connector is 11.5 mm × 4.5 mm. It keeps the standard 19 pins of types A and C, but the pin assignment is different from both.
Type E; Automotive The Automotive Connection System has a locking tab to keep the cable from vibrating loose and a shell to help prevent moisture and dirt from interfering with the signals.
The HDMI alternate mode lets a user connect the reversible USB-C connector with the HDMI source devices (mobile, tablet, laptop). This cable connects to video display/sink devices using any of the native HDMI connectors. This is an HDMI cable, in this case a USB-C to HDMI cable.
Cables
An HDMI cable is composed of four shielded twisted pairs, with impedance of the order of 100 Ω (±15%), plus seven separate conductors. HDMI cables with Ethernet differ in that three of the separate conductors instead form an additional shielded twisted pair (with the CEC/DDC ground as a shield).
Although no maximum length for an HDMI cable is specified, signal attenuation (dependent on the cable's construction quality and conducting materials) limits usable lengths in practice and certification is difficult to achieve for lengths beyond 13 m. HDMI 1.3 defines two cable categories: Category 1-certified cables, which have been tested at 74.25 MHz (which would include resolutions such as 720p60 and 1080i60), and Category 2-certified cables, which have been tested at 340 MHz (which would include resolutions such as 1080p60 and 4K30). Category 1 HDMI cables are marketed as "Standard" and Category 2 HDMI cables as "High Speed". This labeling guideline for HDMI cables went into effect on October 17, 2008. Category 1 and 2 cables can either meet the required parameter specifications for inter-pair skew, far-end crosstalk, attenuation and differential impedance, or they can meet the required non-equalized/equalized eye diagram requirements. A cable of about can be manufactured to Category 1 specifications easily and inexpensively by using 28 AWG (0.081 mm2) conductors. With better quality construction and materials, including 24 AWG (0.205 mm2) conductors, an HDMI cable can reach lengths of up to . Many HDMI cables under 5 meters of length that were made before the HDMI 1.3 specification can work as Category 2 cables, but only Category 2-tested cables are guaranteed to work for Category 2 purposes.
As of the HDMI 1.4 specification, the following cable types are defined for HDMI in general:
Standard HDMI Cableup to 1080i and 720p
Standard HDMI Cable with Ethernet
Standard Automotive HDMI Cable
High Speed HDMI Cable1080p, 4K 30Hz, 3D and deep color
High Speed HDMI Cable with Ethernet
A new certification program was introduced in October 2015 to certify that cables work at the maximum bandwidth of the HDMI 2.0 specification. In addition to expanding the set of cable testing requirements, the certification program introduces an EMI test to ensure cables minimize interference with wireless signals. These cables are marked with an anti-counterfeiting authentication label and are defined as:
Premium High Speed HDMI Cable
Premium High Speed HDMI Cable with Ethernet
In conjunction with the HDMI 2.1 specification, a third category of cable was announced on January 4, 2017, called "48G". Also known as Category 3 HDMI or "Ultra High Speed" HDMI, the cable is designed to support the bandwidth of HDMI 2.1, supporting 4K, 5K, 8K and 10K at 120Hz. The cable is backwards compatible with the earlier HDMI devices, using existing HDMI type A, C and D connectors, and includes HDMI Ethernet.
Ultra High Speed HDMI Cable (48G Cable)4K, 5K, 8K and 10K at 120Hz
Extenders
An HDMI extender is a single device (or pair of devices) powered with an external power source or with the 5V DC from the HDMI source. Long cables can cause instability of HDCP and blinking on the screen, due to the weakened DDC signal that HDCP requires. HDCP DDC signals must be multiplexed with TMDS video signals to comply with HDCP requirements for HDMI extenders based on a single Category 5/Category 6 cable. Several companies offer amplifiers, equalizers and repeaters that can string several standard HDMI cables together. Active HDMI cables use electronics within the cable to boost the signal and allow for HDMI cables of up to ; those based on HDBaseT can extend to 100 meters; HDMI extenders that are based on dual Category 5/Category 6 cable can extend HDMI to ; while HDMI extenders based on optical fiber can extend HDMI to .
Licensing
The HDMI specification is not an open standard; manufacturers need to be licensed by HDMI LA in order to implement HDMI in any product or component. Companies that are licensed by HDMI LA are known as HDMI Adopters.
DVI is the only interface that does not require a license for interfacing HDMI.
HDMI adopters
While earlier versions of HDMI specs are available to the public for download, only adopters have access to the latest standards (HDMI 1.4b/2.1). Only adopters have access to the compliance test specification (CTS) that is used for compliance and certification. Compliance testing is required before any HDMI product can be legally sold.
Adopters have IP rights under Adopter Agreement.
Adopters receive the right to use HDMI logos and TMs on their products and marketing materials.
Adopters are listed on the HDMI website.
Products from adopters are listed and marketed in the official HDMI product finder database.
Adopters receive more exposure through combined marketing, such as the annual HDMI Developers Conference and technology seminars.
HDMI fee structure
There are two annual fee structures associated with being an HDMI adopter:
High-volume (more than 10,000 units) HDMI Adopter Agreement per year.
Low-volume (10,000 units or fewer) HDMI Adopter Agreement plus a flat per unit administration fee.
The annual fee is due upon the execution of the Adopter Agreement, and must be paid on the anniversary of this date each year thereafter.
The royalty fee structure is the same for all volumes. The following variable per-unit royalty is device-based and not dependent on number of ports, chips or connectors:
for each end-user licensed product
if the HDMI logo is used on the product and promotional material, the per-unit fee drops from to .
if HDCP is implemented and HDMI logo is used, the per-unit fee drops from to .
Use of HDMI logo requires compliance testing. Adopters need to license HDCP separately.
The HDMI royalty is only payable on licensed products that will be sold on a stand-alone basis (i.e., that are not incorporated into another licensed product that is subject to an HDMI royalty). For example, if a cable or IC is sold to an adopter who then includes it in a television subject to a royalty, then the cable or IC maker would not pay a royalty, and the television manufacturer would pay the royalty on the final product. If the cable is sold directly to consumers, then the cable would be subject to a royalty.
Versions
HDMI devices and cables are designed based on the HDMI Specification, a document published by HDMI Licensing (through version 1.4b) or the HDMI Forum (from version 2.0 onward). The HDMI Specification defines the minimum baseline requirements that all HDMI devices must adhere to for interoperability, as well as a large set of optional features that HDMI devices may support. The specification is periodically updated to add clarifications or define new capabilities that HDMI devices may implement. Each new version of the specification expands the list of possible features, but does not mandate support for new features in all devices or establish any "classes" of HDMI products which must support certain capabilities. Version numbers do not refer to classes or tiers of products with certain levels of feature support, and as such, HDMI specification "version numbers" are not a method of describing support for specific features or describing the capabilities of an HDMI device or cable.
In 2009, HDMI Licensing banned the use of "version numbers" in labeling HDMI products. Instead, HDMI devices should explicitly declare which features and capabilities they support. For HDMI cables, a speed rating system was established since feature support is not dependent on the cable (apart from inline Ethernet and ARC); the cable only affects the maximum possible speed of the connection. HDMI cables should be labeled with the appropriate speed certification (i.e. Standard Speed, High Speed, or Ultra High Speed), not a "version number".
Version 1.0
HDMI 1.0 was released on December 9, 2002, and is a single-cable digital audio/video connector interface. The link architecture is based on DVI, using exactly the same video transmission format but sending audio and other auxiliary data during the blanking intervals of the video stream. HDMI 1.0 allows a maximum TMDS clock of 165MHz ( bandwidth per link), the same as DVI. It defines two connectors called type A and type B, with pinouts based on the Single-Link DVI-D and Dual-Link DVI-D connectors respectively, though the type B connector was never used in any commercial products. HDMI 1.0 uses TMDS encoding for video transmission, giving it of video bandwidth ( or at 60Hz) and 8-channel LPCM/192 kHz/24-bit audio. HDMI 1.0 requires support for RGB video, with optional support for 4:4:4 and 4:2:2 (mandatory if the device has support for on other interfaces). Color depth of 10bpc (30bit/px) or 12bpc (36bit/px) is allowed when using 4:2:2 subsampling, but only 8bpc (24bit/px) color depth is permitted when using RGB or 4:4:4. Only the Rec. 601 and Rec. 709 color spaces are supported. HDMI 1.0 allows only specific pre-defined video formats, including all the formats defined in EIA/CEA-861-B and some additional formats listed in the HDMI Specification itself. All HDMI sources/sinks must also be capable of sending/receiving native Single-Link DVI video and be fully compliant with the DVI Specification.
Version 1.1
HDMI 1.1 was released on May 20, 2004, and added support for DVD-Audio.
Version 1.2
HDMI 1.2 was released on August 8, 2005, and added the option of One Bit Audio, used on Super Audio CDs, at up to 8 channels. To make HDMI more suitable for use on PC devices, version 1.2 also removed the requirement that only explicitly supported formats be used. It added the ability for manufacturers to create vendor-specific formats, allowing any arbitrary resolution and refresh rate rather than being limited to a pre-defined list of supported formats. In addition, it added explicit support for several new formats including 720p at 100 and 120 Hz and relaxed the pixel format support requirements so that sources with only native RGB output (PC sources) would not be required to support output.
HDMI 1.2a was released on December 14, 2005 and fully specifies Consumer Electronic Control (CEC) features, command sets and CEC compliance tests.
Version 1.3
HDMI 1.3 was released on June 22, 2006, and increased the maximum TMDS clock to 340MHz (). Like previous versions, it uses TMDS encoding, giving it a maximum video bandwidth of (sufficient for at 144Hz or at 75Hz). It added support for 10bpc, 12bpc, and 16bpc color depth (30, 36, and 48bit/px), called deep color. It also added support for the xvYCC color space, in addition to the ITU-R BT.601 and BT.709 color spaces supported by previous versions, and added the ability to carry metadata defining color gamut boundaries. It also optionally allows output of Dolby TrueHD and DTS-HD Master Audio streams for external decoding by AV receivers. It incorporates automatic audio syncing (audio video sync) capability. It defined cable Categories 1 and 2, with Category 1 cable being tested up to 74.25MHz and Category 2 being tested up to 340 MHz. It also added the new HDMI type C "Mini" connector for portable devices.
HDMI 1.3a was released on November 10, 2006, and had cable and sink modifications for HDMI type C, source termination recommendations, and removed undershoot and maximum rise/fall time limits. It also changed CEC capacitance limits, and CEC commands for timer control were brought back in an altered form, with audio control commands added. It also added the optional ability to stream SACD in its bitstream DST format rather than uncompressed raw DSD. HDMI 1.3a is available to download free of charge, after registration.
Version 1.4
HDMI 1.4 was released on June 5, 2009, and first came to market after Q2 of 2009. Retaining the bandwidth of the previous version, HDMI 1.4 defined standardized timings to use for 40962160 at 24Hz, 38402160 at 24, 25, and 30Hz, and added explicit support for 19201080 at 120Hz with CTA-861 timings. It also added an HDMI Ethernet Channel (HEC) that accommodates a Ethernet connection between the two HDMI connected devices so they can share an Internet connection, introduced an audio return channel (ARC), 3D Over HDMI, a new Micro HDMI Connector, an expanded set of color spaces with the addition of sYCC601, Adobe RGB and Adobe YCC601, and an Automotive Connection System. HDMI 1.4 defined several stereoscopic 3D formats including field alternative (interlaced), frame packing (a full resolution top-bottom format), line alternative full, side-by-side half, side-by-side full, 2D + depth, and 2D + depth + graphics + graphics depth (WOWvx). HDMI 1.4 requires that 3D displays implement the frame packing 3D format at either 720p50 and 1080p24 or 720p60 and 1080p24. High Speed HDMI cables as defined in HDMI 1.3 work with all HDMI 1.4 features except for the HDMI Ethernet Channel, which requires the new High Speed HDMI Cable with Ethernet defined in HDMI 1.4.
HDMI 1.4a was released on March 4, 2010, and added two mandatory 3D formats for broadcast content, which was deferred with HDMI 1.4 pending the direction of the 3D broadcast market. HDMI 1.4a has defined mandatory 3D formats for broadcast, game, and movie content. HDMI 1.4a requires that 3D displays implement the frame packing 3D format at either 720p50 and 1080p24 or 720p60 and 1080p24, side-by-side horizontal at either 1080i50 or 1080i60, and top-and-bottom at either 720p50 and 1080p24 or 720p60 and 1080p24.
HDMI 1.4b was released on October 11, 2011, containing only minor clarifications to the 1.4a document. HDMI 1.4b is the last version of the standard that HDMI LA is responsible for. All later versions of the HDMI Specification are produced by the HDMI Forum, created on October 25, 2011.
Version 2.0
HDMI 2.0, referred to by some manufacturers as HDMI UHD, was released on September 4, 2013.
HDMI 2.0 increases the maximum bandwidth to . HDMI 2.0 uses TMDS encoding for video transmission like previous versions, giving it a maximum video bandwidth of . This enables HDMI 2.0 to carry 4K video at 60 Hz with 24 bit/px color depth. Other features of HDMI 2.0 include support for the Rec. 2020 color space, up to 32 audio channels, up to 1536 kHz audio sample frequency, dual video streams to multiple users on the same screen, up to four audio streams, 4:2:0 chroma subsampling, 25 fps 3D formats, support for the 21:9 aspect ratio, dynamic synchronization of video and audio streams, the HE-AAC and DRA audio standards, improved 3D capability, and additional CEC functions.
HDMI 2.0a was released on April 8, 2015, and added support for High Dynamic Range (HDR) video with static metadata.
HDMI 2.0b was released March 2016. HDMI 2.0b initially supported the same HDR10 standard as HDMI 2.0a as specified in the CTA-861.3 specification. In December 2016 additional support for HDR Video transport was added to HDMI 2.0b in the CTA-861-G specification, which extends the static metadata signaling to include hybrid log–gamma (HLG).
Version 2.1
HDMI 2.1 was officially announced by the HDMI Forum on January4, 2017, and was released on November 28, 2017. It adds support for higher resolutions and higher refresh rates, including 4K 120Hz and 8K 60Hz. HDMI 2.1 also introduces a new HDMI cable category called Ultra High Speed (referred to as 48G during development), which certifies cables at the new higher speeds that these formats require. Ultra High Speed HDMI cables are backwards compatible with older HDMI devices, and older cables are compatible with new HDMI 2.1 devices, though the full bandwidth is only supported with the new cables.
Some systems may not be able to use HDMI 2.1 because the HDMI Forum is preventing its use in open source implementations (such as Linux open source drivers). Users of those systems may need to use DisplayPort instead to access high resolutions and speeds.
The following features were added to the HDMI 2.1 Specification:
Maximum supported format is 10K at 120Hz
Dynamic HDR for specifying HDR metadata on a scene-by-scene or even a frame-by-frame basis
Note: While HDMI 2.1 did standardize transport of dynamic HDR metadata over HDMI, in actuality it only formalized dynamic metadata interfaces already utilized by Dolby Vision and HDR10+ in HDMI 2.0, which is why neither Dolby Vision nor HDR10+ require HDMI 2.1 to function properly.
Display Stream Compression (DSC) 1.2 is used for video formats higher than 8K with 4:2:0 chroma subsampling
High Frame Rate (HFR) for 4K, 8K, and 10K, which adds support for refresh rates up to 120Hz
Enhanced Audio Return Channel (eARC) for object-based audio formats such as Dolby Atmos and DTS:X
Enhanced refresh rate and latency reduction features:
Variable Refresh Rate (VRR) reduces or eliminates lag, stutter and frame tearing for more fluid motion in games
Quick Media Switching (QMS) for movies and video eliminates the delay that can result in blank screens before content begins to be displayed
Quick Frame Transport (QFT) reduces latency by bursting individual pictures across the HDMI link as fast as possible when the link's hardware supports more bandwidth than the minimum amount needed for the resolution and frame rate of the content. With QFT, individual pictures arrive earlier and some hardware blocks can be fully powered off for longer periods of time between pictures to reduce heat generation and extend battery life.
Auto Low Latency Mode (ALLM)When a display device supports the option to either optimize its pixel processing for best latency or best pixel processing, ALLM allows the current HDMI source device to automatically select, based on its better understanding of the nature of its own content, which mode the user would most likely prefer.
Video formats that require more bandwidth than (4K 60Hz 8bpc RGB), such as 4K 60Hz 10bpc (HDR), 4K 120Hz, and 8K 60Hz, may require the new "Ultra High Speed" or "Ultra High Speed with Ethernet" cables. HDMI 2.1's other new features are supported with existing HDMI cables.
The increase in maximum bandwidth is achieved by increasing both the bitrate of the data channels and the number of channels. Previous HDMI versions use three data channels (each operating at up to in HDMI 2.0, or up to in HDMI 1.4), with an additional channel for the TMDS clock signal, which runs at a fraction of the data channel speed (one tenth the speed, or up to 340MHz, for signaling rates up to ; one fortieth the speed, or up to 150MHz, for signaling rates between 3.4 and ). HDMI 2.1 doubles the signaling rate of the data channels to . The structure of the data has been changed to use a new packet-based format with an embedded clock signal, which allows what was formerly the TMDS clock channel to be used as a fourth data channel instead, increasing the signaling rate across that channel to 12Gbit/s as well. These changes increase the aggregate bandwidth from (3 × ) to (4 × ), a 2.66× improvement in bandwidth. In addition, the data is transmitted more efficiently by using a 16b/18b encoding scheme, which uses a larger percentage of the bandwidth for data rather than DC balancing compared to the TMDS scheme used by previous versions (88.% compared to 80%). This, in combination with the 2.66× bandwidth, raises the maximum data rate of HDMI 2.1 from to 42.Gbit/s. Subtracting overhead for FEC, the usable data rate is approximately , around 2.92× the data rate of HDMI 2.0.
The bandwidth provided by HDMI 2.1 is enough for 8K resolution at approximately 50Hz, with 8bpc RGB or 4:4:4 color. To achieve even higher formats, HDMI 2.1 can use Display Stream Compression (DSC) with a compression ratio of up to . Using DSC, formats up to 8K () 120Hz or 10K () 100Hz at 8bpc RGB/4:4:4 are possible. Using with 4:2:2 or 4:2:0 chroma subsampling in combination with DSC can allow for even higher formats.
HDMI 2.1a was released on February 15, 2022, and added support for Source-Based Tone Mapping (SBTM).
HDMI 2.1b was released on August 10, 2023.
Version 2.2
HDMI 2.2 was announced on January 6, 2025. It will be released in the first half of 2025. The maximum allowed bit rate is increased to .
Version comparison
Main specifications
Refresh frequency limits for common resolutions
The maximum limits for TMDS transmission are calculated using standard data rate calculations. For FRL transmission, the limits are calculated using the capacity computation algorithm provided by the HDMI Specification. All calculations assume uncompressed RGB video with CVT-RB v2 timing. Maximum limits may differ if compression (i.e. DSC) or 4:2:0 chroma subsampling are used.
Display manufacturers may also use non-standard blanking intervals (a Vendor-Specific Timing Format as defined in the HDMI Specification) rather than CVT-RB v2 to achieve even higher frequencies when bandwidth is a constraint. The refresh frequencies in the below table do not represent the absolute maximum limit of each interface, but rather an estimate based on a modern standardized timing formula. The minimum blanking intervals (and therefore the exact maximum frequency that can be achieved) will depend on the display and how many secondary data packets it requires, and therefore will differ from model to model.
Refresh frequency limits for standard video
HDMI 1.0 and 1.1 are restricted to transmitting only certain video formats, defined in EIA/CEA-861-B and in the HDMI Specification itself. HDMI 1.2 and all later versions allow any arbitrary resolution and frame rate (within the bandwidth limit). Formats that are not supported by the HDMI Specification (i.e., no standardized timings defined) may be implemented as a vendor-specific format. Successive versions of the HDMI Specification continue to add support for additional formats (such as 4K resolutions), but the added support is to establish standardized timings to ensure interoperability between products, not to establish which formats are or are not permitted. Video formats do not require explicit support from the HDMI Specification in order to be transmitted and displayed.
Individual products may have heavier limitations than those listed below, since HDMI devices are not required to support the maximum bandwidth of the HDMI version that they implement. Therefore, it is not guaranteed that a display will support the refresh rates listed in this table, even if the display has the required HDMI version.
Uncompressed 8bpc (24bit/px) color depth and RGB or 4:4:4 color format are assumed on this table except where noted.
Refresh frequency limits for HDR10 video
HDR10 requires 10bpc (30bit/px) color depth, which uses 25% more bandwidth than standard 8bpc video.
Uncompressed 10bpc color depth and RGB or 4:4:4 color format are assumed on this table except where noted.
Feature support
The features defined in the HDMI Specification that an HDMI device may implement are listed below. For historical interest, the version of the HDMI Specification in which the feature was first added is also listed. All features of the HDMI Specification are optional; HDMI devices may implement any combination of these features.
Although the "HDMI version numbers" are commonly misused as a way of indicating that a device supports certain features, this notation has no official meaning and is considered improper by HDMI Licensing. There is no officially-defined correlation between features supported by a device and any claimed "version numbers", as version numbers refer to historical editions of the HDMI Specification document, not to particular classes of HDMI devices. Manufacturers are forbidden from describing their devices using HDMI version numbers, and are required to identify support for features by listing explicit support for them, but the HDMI forum has received criticism for lack of enforcement of these policies.
Full HD Blu-ray Disc and HD DVD video (version 1.0)
Consumer Electronic Control (CEC) (version 1.0)
DVD-Audio (version 1.1)
Super Audio CD (DSD) (version 1.2)
Auto Lip-Sync Correction (version 1.3)
Dolby TrueHD / DTS-HD Master Audio bitstream capable (version 1.3)
Updated list of CEC commands (version 1.3a)
3D video (version 1.4)
Ethernet channel (100Mbit/s) (version 1.4)
Audio return channel (ARC) (version 1.4)
4 audio streams (version 2.0)
Dual View (version 2.0)
Perceptual quantizer HDR EOTF (SMPTE ST 2084) (version 2.0a)
Hybrid log–gamma (HLG) HDR EOTF (version 2.0a)
Static HDR metadata (SMPTE ST 2086) (version 2.0a)
Dynamic HDR metadata (SMPTE ST 2094) (version 2.0b)
Enhanced audio return channel (eARC) (version 2.1)
Variable Refresh Rate (VRR) (version 2.1)
Quick Media Switching (QMS) (version 2.1)
Quick Frame Transport (QFT) (version 2.1)
Auto Low Latency Mode (ALLM) (version 2.1)
Display Stream Compression (DSC) (version 2.1)
Source-Based Tone Mapping (SBTM) (version 2.1a)
Display Stream Compression
Display Stream Compression (DSC) is a VESA-developed video compression algorithm designed to enable increased display resolutions and frame rates over existing physical interfaces, and make devices smaller and lighter, with longer battery life.
Applications
Blu-ray Disc and HD DVD players
Blu-ray Disc and HD DVD, introduced in 2006, offer high-fidelity audio features that require HDMI for best results. HDMI 1.3 can transport Dolby Digital Plus, Dolby TrueHD, and DTS-HD Master Audio bitstreams in compressed form. This capability allows for an AV receiver with the necessary decoder to decode the compressed audio stream. The Blu-ray specification does not include video encoded with either deep color or xvYCC; thus, HDMI 1.0 can transfer Blu-ray discs at full video quality.
The HDMI 1.4 specification (released in 2009) added support for 3D video and is used by all Blu-ray 3D compatible players.
The Blu-ray Disc Association (BDA) spokespersons have stated (Sept. 2014 at IFA show in Berlin, Germany) that the Blu-ray, Ultra HD players, and 4K discs are expected to be available starting in the second half to 2015. It is anticipated that such Blu-ray UHD players will be required to include a HDMI 2.0 output that supports HDCP 2.2.
Blu-ray permits secondary audio decoding, whereby the disc content can tell the player to mix multiple audio sources together before final output. Some Blu-ray and HD DVD players can decode all of the audio codecs internally and can output LPCM audio over HDMI. Multichannel LPCM can be transported over an HDMI connection, and as long as the AV receiver implements multichannel LPCM audio over HDMI and implements HDCP, the audio reproduction is equal in resolution to HDMI 1.3 bitstream output. Some low-cost AV receivers, such as the Onkyo TX-SR506, do not allow audio processing over HDMI and are labelled as "HDMI pass through" devices. Virtually all modern AV Receivers now offer HDMI 1.4 inputs and outputs with processing for all of the audio formats offered by Blu-ray Discs and other HD video sources. During 2014 several manufacturers introduced premium AV Receivers that include one, or multiple, HDMI 2.0 inputs along with a HDMI 2.0 output(s). However, not until 2015 did most major manufacturers of AV receivers also support HDCP 2.2 as needed to support certain high quality UHD video sources, such as Blu-ray UHD players.
Digital cameras and camcorders
Most consumer camcorders, as well as many digital cameras, are equipped with a mini-HDMI connector (type C connector).
Some cameras also have 4K capability,
although cameras capable of HD video often include an HDMI interface for playback or even live preview, the image processor and the video processor of cameras usable for uncompressed video must be able to deliver the full image resolution at the specified frame rate in real time without any missing frames causing jitter. Therefore, usable uncompressed video out of HDMI is often called "clean HDMI".
Personal computers
Personal computer (PCs) with a DVI interface are capable of video output to an HDMI-enabled monitor. Some PCs include an HDMI interface and may also be capable of HDMI audio output, depending on specific hardware. For example, Intel's motherboard chipsets since the 945G and NVIDIA's GeForce 8200/8300 motherboard chipsets are capable of 8-channel LPCM output over HDMI. Eight-channel LPCM audio output over HDMI with a video card was first seen with the ATI Radeon HD 4850, which was released in June 2008 and is implemented by other video cards in the ATI Radeon HD 4000 series. Linux can drive 8-channel LPCM audio over HDMI if the video card has the necessary hardware and implements the Advanced Linux Sound Architecture (ALSA). The ATI Radeon HD 4000 series implements ALSA. Cyberlink announced in June 2008 that they would update their PowerDVD playback software to allow 192 kHz/24-bit Blu-ray Disc audio decoding in Q3-Q4 of 2008. Corel's WinDVD 9 Plus currently has 96 kHz/24-bit Blu-ray Disc audio decoding.
Even with an HDMI output, a computer may not be able to produce signals that implement HDCP, Microsoft's Protected Video Path, or Microsoft's Protected Audio Path. Several early graphic cards were labelled as "HDCP-enabled" but did not have the hardware needed for HDCP; this included some graphic cards based on the ATI X1600 chipset and certain models of the NVIDIA Geforce 7900 series. The first computer monitors that could process HDCP were released in 2005; by February 2006 a dozen different models had been released. The Protected Video Path was enabled in graphic cards that had HDCP capability, since it was required for output of Blu-ray Disc and HD DVD video. In comparison, the Protected Audio Path was required only if a lossless audio bitstream (such as Dolby TrueHD or DTS-HD MA) was output. Uncompressed LPCM audio, however, does not require a Protected Audio Path, and software programs such as PowerDVD and WinDVD can decode Dolby TrueHD and DTS-HD MA and output it as LPCM. A limitation is that if the computer does not implement a Protected Audio Path, the audio must be downsampled to 16-bit 48 kHz but can still output at up to 8 channels. No graphic cards were released in 2008 that implemented the Protected Audio Path.
The Asus Xonar HDAV1.3 became the first HDMI sound card that implemented the Protected Audio Path and could both bitstream and decode lossless audio (Dolby TrueHD and DTS-HD MA), although bitstreaming is only available if using the ArcSoft TotalMedia Theatre software. It has an HDMI 1.3 input/output, and Asus says that it can work with most video cards on the market.
In September 2009, AMD announced the ATI Radeon HD 5000 series video cards, which have HDMI 1.3 output (deep color, xvYCC wide gamut capability and high bit rate audio), 8-channel LPCM over HDMI, and an integrated HD audio controller with a Protected Audio Path that allows bitstream output over HDMI for AAC, Dolby AC-3, Dolby TrueHD and DTS-HD Master Audio formats. The ATI Radeon HD 5870 released in September 2009 is the first video card that allows bitstream output over HDMI for Dolby TrueHD and DTS-HD Master Audio. The AMD Radeon HD 6000 series implements HDMI 1.4a. The AMD Radeon HD 7000 series implements HDMI 1.4b.
In December 2010, it was announced that several computer vendors and display makers including Intel, AMD, Dell, Lenovo, Samsung, and LG would stop using LVDS (actually, FPD-Link) from 2013 and legacy DVI and VGA connectors from 2015, replacing them with DisplayPort and HDMI.
On August 27, 2012, Asus announced a new monitor that produces its native resolution of 2560×1440 via HDMI 1.4.
On September 18, 2014, Nvidia launched GeForce GTX 980 and GTX 970 (with GM204 chip) with HDMI 2.0 support. On January 22, 2015, GeForce GTX 960 (with GM206 chip) launched with HDMI 2.0 support. On March 17, 2015, GeForce GTX TITAN X (GM200) launched with HDMI 2.0 support. On June 1, 2015, GeForce GTX 980 Ti (with GM200 chip) launched with HDMI 2.0 support. On August 20, 2015, GeForce GTX 950 (with GM206 chip) launched with HDMI 2.0 support.
On May 6, 2016, Nvidia launched the GeForce GTX 1080 (GP104 GPU) with HDMI 2.0b support.
On September 1, 2020, Nvidia launched the GeForce RTX 30 series, the world's first discrete graphics cards with support for the full bandwidth with Display Stream Compression 1.2 of HDMI 2.1.
Gaming consoles
Beginning with the seventh generation of video game consoles, most consoles support HDMI. Video game consoles that support HDMI include the Xbox 360 (except most pre-2007 models) (1.2a), Xbox One (1.4b), Xbox One S (2.0a), Xbox One X (2.0b), PlayStation 3 (1.3a), PlayStation 4 (1.4b), PlayStation 4 Pro (2.0a), Wii U (1.4a), Nintendo Switch (1.4b), Nintendo Switch (OLED model) (2.0a), Xbox Series X and Series S (2.1), and PlayStation 5 (2.1).
Tablet computers
Some tablet computers implement HDMI using Micro-HDMI (type D) port, while others like the Eee Pad Transformer implement the standard using mini-HDMI (type C) ports. All iPad models have a special A/V adapter that converts Apple's Lightning connector to a standard HDMI (type A) port. Samsung has a similar proprietary thirty-pin port for their Galaxy Tab 10.1 that could adapt to HDMI as well as USB drives. The Dell Streak 5 smartphone/tablet hybrid is capable of outputting over HDMI. While the Streak uses a PDMI port, a separate cradle adds HDMI compatibility. Some tablets running Android OS provide HDMI output using a mini-HDMI (type C) port. Most new laptops and desktops now have built in HDMI as well.
Mobile phones
Many mobile phones can produce an output of HDMI video via a micro-HDMI connector, SlimPort, MHL or other adapter.
Legacy compatibility
HDMI can only be used with older analog-only devices (using connections such as SCART, VGA, RCA, etc.) by means of a digital-to-analog converter or AV receiver, as the interface does not carry any analog signals (unlike DVI, where devices with DVI-I ports accept or provide either digital or analog signals). Cables are available that contain the necessary electronics, but it is important to distinguish these active converter cables from passive HDMI to VGA cables (which are typically cheaper as they don't include any electronics). The passive cables are only useful if a user has a device that is generating or expecting HDMI signals on a VGA connector, or VGA signals on an HDMI connector; this is a non-standard feature, not implemented by most devices.
HDMI Alternate Mode for USB type-C
The HDMI Alternate Mode for USB-C allows HDMI-enabled sources with a USB-C connector to directly connect to standard HDMI display devices, without requiring an adapter. The standard was released in September 2016, and supports all HDMI 1.4b features such as video resolutions up to Ultra HD 30 Hz and CEC. Previously, the similar DisplayPort Alternate Mode could be used to connect to HDMI displays from USB type-C sources, but where in that case active adapters were required to convert from DisplayPort to HDMI, HDMI Alternate Mode connects to the display natively.
The Alternate Mode reconfigures the four SuperSpeed differential pairs present in USB-C to carry the three HDMI TMDS channels and the clock signal. The two Sideband Use pins (SBU1 and SBU2) are used to carry the HDMI Ethernet and Audio Return Channel and the Hot Plug Detect functionality (HEAC+/Utility pin and HEAC−/HPD pin). As there are not enough reconfigurable pins remaining in USB-C to accommodate the DDC clock (SCL), DDC data (SDA), and CECthese three signals are bridged between the HDMI source and sink via the USB Power Delivery 2.0 (USB-PD) protocol, and are carried over the USB-C Configuration Channel (CC) wire. This is possible because the cable is electronically marked (i.e., it contains a USB-PD node) that serves to tunnel the DDC and CEC from the source over the Configuration Channel to the node in the cable, these USB-PD messages are received and relayed to the HDMI sink as regenerated DDC (SCL and SDA signals), or CEC signals.
As stated at CES in January 2023, HDMI Alternate Mode for USB Type-C is no longer being updated as there are no known products using this protocol, reducing its relevance in the current market. This will reduce consumer confusion as DisplayPort Alternate Mode is the primary video protocol of choice over USB-C.
Relationship with DisplayPort
The DisplayPort audio/video interface was introduced in May 2006. Historically, HDMI Licensing LLC was publicly dismissive of DisplayPort's position in the industry, with its president stating in a 2009 interview that "there are certainly some PCs that have
DisplayPort connectors on them, but these are niche applications that have not taken hold in the market."
In recent years, DisplayPort connectors have become a common feature of premium products—displays, desktop computers, and video cards; most of the companies producing DisplayPort equipment are in the computer sector. The DisplayPort website states that DisplayPort is expected to complement HDMI, but 100% of HD and UHD TVs had HDMI connectivity. DisplayPort supported some advanced features which are useful for multimedia content creators and gamers (e.g., 5K, Adaptive-Sync), which was the reason most GPUs have DisplayPort. These features were added to the official HDMI specification slightly later, but with the introduction of HDMI 2.1, these gaps are already leveled off (e.g., VRR / Variable Refresh Rate).
DisplayPort uses a self-clocking, micro-packet-based protocol that allows for a variable number of differential pair lanes as well as flexible allocation of bandwidth between audio and video, and allows encapsulating multi-channel compressed audio formats in the audio stream. DisplayPort 1.2 supports multiple audio/video streams, variable refresh rate (FreeSync), and Dual-mode transmitters compatible with HDMI 1.2 or 1.4. Revision 1.3 increases overall transmission bandwidth to with the new HBR3 mode featuring per lane; it requires Dual-mode with mandatory HDMI 2.0 compatibility and HDCP 2.2. Revision 1.4 added Display Stream Compression (DSC), support for the BT.2020 color space, and HDR10 extensions from CTA-861.3, including static and dynamic metadata. Revision 1.4a was published in April 2018, updating DisplayPort's DSC implementation from 1.2 to 1.2a. Revision 2.0 increased overall bandwidth from 25.92 to , enabling increased resolutions and refresh rates, increasing the resolutions and refresh rates with HDR support, and other related improvements. Revision 2.1 was published in October 2022, incorporating the new DP40 and DP80 cable certifications, which require proper operation at the UHBR10 (40Gbit/s) and UHBR20 (80Gbit/s) speeds introduced in version 2.0, and a bandwidth management feature to enable DisplayPort tunnelling to coexist with other I/O data traffic more efficiently over a USB4/USB Type-C connection.
The DisplayPort features an adapter detection mechanism enabling dual-mode operation and the transmission of TMDS signals allowing the conversion to DVI and HDMI 1.2/1.4/2.0 signals using a passive adapter. The same external connector is used for both protocolswhen a DVI/HDMI passive adapter is attached, the transmitter circuit switches to TMDS mode. DisplayPort Dual-mode ports and cables/adapters are typically marked with the DisplayPort++ logo. Thunderbolt ports with mDP connector also supports Dual-mode passive HDMI adapters/cables. Conversion to dual-link DVI and component video (VGA/YPbPr) requires active powered adapters.
The USB 3.1 type-C connector is increasingly the standard video connector, replacing legacy video connectors such as mDP, Thunderbolt, HDMI, and VGA in mobile devices. USB-C connectors can transmit DisplayPort video to docks and displays using standard USB type-C cables or type-C to DisplayPort cables and adapters; USB-C also supports HDMI adapters that actively convert from DisplayPort to HDMI 1.4 or 2.0. DisplayPort Alternate Mode for USB type-C specification was published in 2015. USB type-C chipsets are not required to include Dual-mode, so passive DP-HDMI adapters do not work with type-C sources. A specification for "HDMI Alternate Mode for USB type-C" was released in 2016, but was discontinued in 2023, with HDMI Licensing Administration stating they knew of no adapter having ever been produced.
DisplayPort is royalty-free, though patent pool administrator Via LA attempts to collect a $0.20 per-device charge for a bulk license to patents it regards as essential to the DisplayPort specification, while HDMI has an annual fee of US$10,000 and a per unit royalty rate of between $0.04 and $0.15.
HDMI has had a few advantages over DisplayPort, such as ability to carry Consumer Electronics Control (CEC) signals since its first generation (DisplayPort 1.3, introduced in 2014, is the earliest DisplayPort generation which can carry CEC signals).
Relationship with MHL
Mobile High-Definition Link (MHL) is an adaptation of HDMI intended to connect mobile devices such as smartphones and tablets to high-definition televisions (HDTVs) and displays. Unlike DVI, which is compatible with HDMI using only passive cables and adapters, MHL requires that the HDMI socket be MHL-enabled, otherwise an active adapter (or dongle) is required to convert the signal to HDMI. MHL is developed by a consortium of five consumer electronics manufacturers, several of which are also behind HDMI.
MHL pares down the three TMDS channels in a standard HDMI connection to a single one running over any connector that provides at least five pins. This lets existing connectors in mobile devicessuch as micro-USBbe used, avoiding the need for additional dedicated video output sockets. The USB port switches to MHL mode when it detects a compatible device is connected.
In addition to the features in common with HDMI (such as HDCP encrypted uncompressed high-definition video and eight-channel surround sound), MHL also adds the provision of power charging for the mobile device while in use, and also enables the TV remote to control it. Although support for these additional features requires connection to an MHL-enabled HDMI port, power charging can also be provided when using active MHL to HDMI adapters (connected to standard HDMI ports), provided there is a separate power connection to the adapter.
Like HDMI, MHL defines a USB-C Alternate Mode to support the MHL standard over USB-C connections.
Version 1.0 supported 720p/1080i 60 Hz (RGB/4:4:4 pixel encoding) with a bandwidth of . Versions 1.3 and 2.0 added support for 1080p 60 Hz ( 4:2:2) with a bandwidth of in PackedPixel mode. Version 3.0 increased the bandwidth to to support Ultra HD (3840 × 2160) 30 Hz video, and also changed from being frame-based, like HDMI, to packet-based.
The fourth version, superMHL, increased bandwidth by operating over multiple TMDS differential pairs (up to a total of six) allowing a maximum of . The six lanes are supported over a reversible 32-pin superMHL connector, while four lanes are supported over USB-C Alternate Mode (only a single lane is supported over micro-USB/HDMI). Display Stream Compression (DSC) is used to allow up to 8K Ultra HD (7680 × 4320) 120 Hz HDR video, and to support Ultra HD 60 Hz video over a single lane.
| Technology | User interface | null |
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