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1991704 | https://en.wikipedia.org/wiki/Nymphaea%20nouchali%20var.%20caerulea | Nymphaea nouchali var. caerulea | Nymphaea nouchali var. caerulea, is a water lily in the genus Nymphaea, a botanical variety of Nymphaea nouchali.
It is an aquatic plant of freshwater lakes, pools and rivers, naturally found throughout most of the eastern half of Africa, as well as parts of southern Arabia, but has also been spread to other regions as an ornamental plant. It was grown by the Ancient Egyptian civilization and had significance in their religion.
It can tolerate the roots being in anoxic mud in nutritionally poor conditions, and can become a dominant plant in deeper water in such habitats. It is associated with a species of snail, which is one of the main hosts of the pathogen causing human schistosomiasis. The underwater rhizomes are edible.
Nymphaea caerulea, first described by Marie Jules César Savigny in 1798, was later classified as a variety of Nymphaea nouchali by Bernard Verdcourt in 1989. Though it is still most commonly referred to as a variety of Nymphaea nouchali, recent phylogenetic studies have problematized the taxonomy. When defined taxonomically as Nymphaea nouchali var. caerulea, it is considered synonymous with Nymphaea capensis.
Description
Vegetative characteristics
This is an aquatic (euhydrophyte) herb with a tuberous rhizome. That is to say, it has small tubers that may develop into short vertical rhizomes. It is a perennial. One plant can spread over an area of about 1 metre.
The peltate leaves have long petioles and have leaf blades (lamina) which are by cm in size. The leaves are polymorphic, changing in form and texture depending if they are underwater or floating. These laminae have a chartaceous texture and can be glabrous or densely covered in pubescent hairs. The shape is incised-cordate and orbicular or subelliptic, with an acute or caudate apex. The two lobes can overlap somewhat or be slightly apart from each other. The upper surface of the lamina is smooth, but the underside has conspicuously raised, green or rarely reddish or reddish-purple veins. There are eight to eleven primary lateral veins on each side of the midrib. There are six to eight pairs of secondary veins arising from the midrib. The primary veins form a pattern of closed, elongated areas stretching to more than two thirds of the way to the margin of the leaf. The leaf margin is entire towards the apex or more-or-less irregularly sinuate-lobulate throughout its entirety. The petioles are thick, blackish green and spongy. They continue to lengthen as they age, pushing older leaves towards the margins of the plant.
Generative characteristics
The flowers can be blue, white, mauve or pinkish in colour, but are usually have pale bluish-white to sky-blue or mauve petals, smoothly changing to a pale yellow in the centre of the flower, and are in diameter. There are four sepals; these are coloured green and sometimes purple at the margins, and are by in size. There are 14–20 petals, of which the outermost are as long as the sepals. Their shape is oblong, and their apexes end in blunt or subacute tips. The stamens are densely congested and very numerous, numbering 100–200 or more. The outermost stamens have long appendages. There are 14–24 carpels, with a very short style. There are also carpellary appendages; these are what is known as 'osmophores', structures which serve to attract pollinators without actually rewarding them, thus by deceit. In this case they are visually attractive for bees and exude an odour mimicking food.
The flower buds rise to the surface over a period of two to three days, and when ready, open during the mid-morning, closing near dusk. This ability is controlled by the sepals; when these are cut off, the flower loses the ability to close. The flowers and buds do not rise above the water in the morning, nor do they submerge at night. The flowers last some four days before they start to wither, closing up each night.
The fruit are berries, 2.2 by 3.2 cm and flattened-round in shape. The seeds are ellipsoid and 1.2 mm long. They are smooth, and have a fleshy, bell-shaped aril.
Chemical composition
In the 1970s and 1980s William Emboden suggested that Nymphaea caerulea contained aporphine alkaloids (though photochemical studies were not performed); more recent photochemical analysis did not find aporohine alkaloids and the chemical composition of Nymphaea caerulea was discovered to include compounds such as 7-hydroxyflavone, 4,7-dihydroxyaurone, and 4’-hydroxyaurone, along with methyl vanillate and cinnamyl alcohol.
Cytology
The chromosome count is n = 14. The genome size is 567.24 Mb.
Taxonomy
Nymphaea spectabilis, a purple form known from cultivation, and N. capensis, found throughout eastern, central and southern Africa, as well as a number of other named taxa, were synonymised to N. nouchali var. caerulea in the 1989 addition to the Flora of Tropical East Africa (FTEA) series, a position which has generally been accepted, although some of the authorities in Bangladesh and in the United States disagree.
In 2012 there was a phylogenetic study where N. caerulea was more related to N. gracilis, an endemic of northern Mexico, than it was to N. nouchali. The evolutionary tree was a consensus of ITS2 and matk. According to this study, N. caerulea should not be considered as a variety of N. nouchali. When genomes from the water lily genus (Nymphaea) were published in the journal Nature in 2020, N. caerulea was cited under that name, not as N. nouchali var. caerulea. Another phylogenetic study from 2021 found N. caerulea (as N. capensis) to be closest related to N. colorata, an east African species.
Nymphaea nouchali is itself a taxonomically challenging species, with a distribution that spans Australia, throughout southern Asia, across Africa to the Western Cape. It has many colour forms (with red-coloured forms generally called N. stellata) and has a long history of cultivation. In Africa, following the 1989 FTEA publication, five different varieties are recognised: var. caerulea, the most widespread, ovalifolia, in parts of tropical Southern Africa, petersiana, the same, zanzibarensis, from tropical southern, central and East Africa, and mutandaensis, which is an endemic of Uganda. One of these taxa, var. petersiana, was found to be quite divergent in the 2012 study. If the 2012 study is to be accepted, this may indicate that the African populations of N. nouchali belong to another species than the Asian and Australian type populations, and should likely be renamed N. caerulea as this name has priority over N. capensis.
Publication
It was first described as Nymphaea caerulea Savigny by Marie Jules César Savigny in 1798. Later, it was included in the species Nymphaea nouchali Burm.f. as the variety Nymphaea nouchali var. caerulea (Savigny) Verdc. published by Bernard Verdcourt in 1989.
Classification
It is classified in the Nymphaea subgenus Brachyceras. This subgenus appears to be phylogenetically sound.
Distribution
The native distribution covers North Africa along the Nile and south throughout central, East and Southern Africa. It is common in this range. The conservation status has not been evaluated by the IUCN, but it is considered a species of 'least concern' by the South African National Biodiversity Institute in their Red List of South African Plants.
On the African continent, it occurs, from north to south, westwards to at least Chad, Congo-Brazzaville, the DRC (only in Katanga?), Angola and Namibia. In South Africa this plant is found in every province, as well as in eSwatini, but it is not native to Lesotho and the Western Cape. It also occurs on islands off the eastern African coast: Zanzibar, Madagascar and the Comoro Islands. It is native to Yemen and Oman (in Dhofar) in the southern Arabian Peninsula but, according to Moshe Agami in a 1980 paper, is thought to have become extinct in the wild in Israel.
It has more recently been spread more widely around the world as an ornamental plant, and introduced populations are now found in Bangladesh, Meghalaya, Kerala and Assam in India, Fiji, Mauritius, North Island in New Zealand, New South Wales and Queensland in Australia, Cook Islands, Costa Rica, and throughout eastern South America (in Brazil and Argentina).
There is an introduced population of blue water-lilies originally from East Africa in the US in the state of Florida. This was first identified as N. zanzibarensis, then as N. capensis var. zanzibarensis, but following the 1989 FTEA publication the taxon was moved to N. nouchali var. zanzibarensis. Nonetheless the 1997 addition to the Flora of North America series decided to retain recognition of the local population under the name N. capensis, and this population continues to be recognised under that name in the US.
The naturalised populations in eastern Australia were also thought to be N. capensis var. zanzibarensis, then later N. caerulea var. zanzibarensis, then in 2011 N. capensis, but the plants in the wild are now thought to be N. caerulea. It is considered an environmental weed in Australia.
Ecology
It has a habitat consisting of rivers, lakes and pools. As of 1921, it has been found at elevations of in South Africa.
Although in cultivation it is said to be quite demanding of nutrients, in the quite nutrient-poor Lake Nabugabo in Uganda it is the dominant aquatic plant species, only being replaced by N. lotus in the eastern tip of the lake, and other aquatic genera where it is more shallow. The dense monospecific stands are associated with an Utricularia sp. and Nymphoides indica in one part of the lake, and with Ceratophyllum demersum in certain other bays. The waterlily stands in this lake are especially poor in invertebrate biodiversity, which may reflect the low levels of dissolved oxygen near the sediments in this habitat. In Lake Bisina, Uganda, N. caerulea is most clearly associated with Utricularia reflexa; this may be due to similar ecological niches, it may just mean the small, rootless, free-moving Utricularia simply get snagged on the petioles, but it may indicate some sort of a commensal relationship, with U. reflexa being shaded by the leaves of N. caerulea. Hydrilla verticillata is another plant which seems to sometimes occur together with the waterlily in this lake, as well as in Lake Bunyonyi.
Pollination is entomophilous. In Kirstenbosch Botanical Gardens, South Africa, the flowers are visited by honey bees. In fact, the carpellary appendages in this type of water-lily appear to have evolved specifically to attract bee species in general. In a way, these waterlilies are parasites of the services of bees, attracting the insects by deceit, without actually rewarding them for their labours. In India plants bloom and fruit from May to October.
The fruit suddenly bursts when ripe, and the scattered seed float away. The seed soon sinks. Seeds often make it to the river's edge or lake shore, and can build up a significant seedbank here. These seeds only germinate when heavy rains flood the banks, and they are submerged under a layer of water. In cultivation, the plants take three to four years to flower from seed. In colder climates, the plants lose their leaves and go dormant during the winter, with the rhizomes remaining alive below the water.
Gomphonema gracile is an epiphytic diatom found on N. caerulea in high elevation Lake Naivasha, Kenya. In Kenya, N. caerulea is positively associated with the freshwater snail Biomphalaria pfeifferi, which is a main host of human schistosomiasis. The edible American crayfish Procambarus clarkii eliminates the mollusc, as well as feeding on the water-lily. The crayfish was first introduced to Kenya in 1966 as a species with which to enhance the local fisheries. In Lake Naivasha, N. caerulea was extremely common until the 1970s, and there is still a seedbank around the shores of the lake. Procambarus clarkii was introduced to the lake in 1970, and now supports an annual harvest of a few thousand kilograms, but it may have been responsible for eliminating not only the water-lily in the main lake by 1983, but all native aquatic plant species in this water body. It is not the only potential culprit; invasive mats of exotic floating vegetation have also taken over the lake, two different commercially fishable fish species have been introduced, and the new fisheries upon these three species could all be responsible, or a combination.
Uses
The rootstock of the blue water lily was collected and eaten in western South Africa around 1800, either raw or in curries, in particular by the Cape Malays and farming communities in the Cape, although this practice has now died out.
It has been suggested that it was used in ancient Egypt for religious rituals and sexual enhancement, due to the purported presence of apomorphine, which is also used today to treat erectile dysfunction.
This water lily has been used to produce perfumes since ancient times; it is also used in aromatherapy. According to a multimodal analytical study, traces of Peganum harmala, and Nymphaea nouchali var. caerulea were identified in an Egyptian ritual Bes-vase, of the 2nd century BCE.
Cultivation
It is grown as an ornamental plant for water gardens in tropical to subtropical regions around the world. It is easy to grow in ponds in any part of Southern Africa, including the highveld, and is hardy to -1 °C.
'Valentina's Pale Blue Eyes' is a registered cultivar of this species from 2018, bred in Italy partially from a clone known as 'Rwanda'.
The Longwood Gardens in Kennett Square, Pennsylvania has had Nymphaea caerulea in their water lilly collection since 1963. with photos of a real Nymphaea caerulea posted on their social media as recent as 2019.
Religion and art
Along with the white lotus, Nymphaea lotus, also native to Egypt, the plant and flower are very frequently depicted in Ancient Egyptian art. They have been depicted in numerous stone carvings and paintings, including the walls of the temple of Karnak, and may be associated with rites pertaining to the afterlife. A number of pharaohs' mummies were covered with the petals of the flower. There are indications it was grown in special farms over 4,000 years ago to produce enough flowers for votive offerings, although it was apparently also simply grown as an ornamental in traditional Egyptian garden ponds. N. caerulea was considered extremely significant in Egyptian mythology, regarded as a symbol of the sun, since the flowers are closed at night and open again in the morning. At Heliopolis, the origin of the world was taught to have been when the sun god Ra emerged from a lotus flower growing in "primordial waters". At night, he was believed to retreat into the flower again. Due to its colour, it was identified, in some beliefs, as having been the original container, in a similar manner to an egg, of Atum, and in similar beliefs Ra, both solar deities. As such, its properties form the origin of the "lotus variant" of the Ogdoad cosmogony. It was also the symbol of the Egyptian deity Nefertem.
Legal issues
Nymphaea caerulea has been illegal in Latvia since November 2009. It is a schedule 1 drug. Possession of up to 1 gram are fined up to 280 euros; for second offences within a year period, criminal charges are applied. Possession of larger quantities can be punished by up to 15 years in prison. The plant was banned in Poland in March 2009. Possession and distribution lead to a criminal charge. N. caerulea has been illegal in Russia since April 2009, along with related products such as Salvia divinorum, Argyreia nervosa, and others.
| Biology and health sciences | Nymphaeales | Plants |
1991929 | https://en.wikipedia.org/wiki/Babirusa | Babirusa | The babirusas, also called deer-pigs (), are a genus, Babyrousa, in the swine family found in the Indonesian islands of Sulawesi, Togian, Sula and Buru. All members of this genus were considered part of a single species until 2002, the babirusa, B. babyrussa, but following that was split into several species. This scientific name is restricted to the Buru babirusa from Buru and Sula, whereas the best-known species, the North Sulawesi babirusa, is named B. celebensis. The males have prominent upwards incurving canine tusks, which pierce the flesh in the snout.
All species of babirusa are listed as threatened by the International Union for Conservation of Nature (IUCN).
Classification
The genus is monotypic within the subfamily Babyrousinae, or alternatively considered to form a tribe, Babyrousini, of the subfamily Suinae. To date, only one fossil skull has been found to suggest a larger ancestor.
All members of the genus were considered part of a single species, the babirusa or pig-deer, B. babyrussa. After they were split into several species, this scientific name is restricted to the Buru babirusa from Buru and the Sula Islands, whereas the best-known species, the north Sulawesi babirusa, is named B. celebensis. The split, which uses the phylogenetic species concept, is based on differences in size, amount of hair on the body and tail-tuft, and measurements of the skull and teeth.
Species
B. babyrussa beruensis was described as an extinct, Pleistocene subspecies from southwestern Sulawesi before babirusas were split into multiple species.
Description
Babirusas are notable for the long upper canines in the males. The upper canines of males emerge vertically from the alveolar process, penetrating through the skin and curving backward over the front of the face and towards the forehead. The lower canines also grow upwards. The canines of females are either reduced or absent. The structure of the male's canines varies by species. In the golden babirusa, the upper canines are short and slender with the alveolar rotated forward to allow the lower canines to cross the lateral view. The Togian babirusa also has the same characteristics and the upper canines always converge. The North Sulawesi babirusa has long and thick upper canines with a vertically implanted alveolar. This causes the upper canines to emerge vertically and not cross with the lower canines.
Babirusas also vary by species in other characteristics. The golden babirusa has a long, thick pelage that is white, creamy gold, black or gold overall, and black at the rump. The pelage of the Togian babirusa is also long but not as that of the golden babirusa. The Togian babirusa has a tawny, brown, or black pelage that is darker on the upper parts than in the lower parts. The North Sulawesi babirusa has very short hair and appears bald. The female babirusa has only one pair of teats.
Biology and ecology
Babirusas are native to Sulawesi, some of the Togian Islands, the Sula Islands, and Buru. In Sulawesi, they range from the Minahasa Peninsula to the provinces of South Sulawesi and Southeast Sulawesi. Although they are present on both Sulawesi and Sula, they are not found on the large islands between the two, the Banggai Archipelago. It has been hypothesized that the unusual distribution may be due to their being transported by humans as gifts bestowed by native royalty. The preferred habitat of babirusa is tropical rainforest along river banks. It appears that they have been confined to the higher grounds in the interior despite occurring in lowland areas near coasts in the past. They are active during the daytime. Like all pig species, babirusa has an omnivorous diet with an intestinal tract similar to that of the domestic pig. The stomach diverticulum of a babirusa is enlarged which may indicate that it is a ruminant but evidence shows otherwise. Because it does not have a rostral bone in the nose, a babirusa does not dig with its snout like other pigs do except in mud and swampy grounds. The diet of the babirusa includes leaves, roots, fruits and animal material. The strong jaws of a babirusa are capable of easily cracking hard nuts.
Males tend to live solitarily while adult females can be found in groups with young. Groups of females and young may number up to 84 individuals, most of which contain no adult males. Males rarely travel in pairs or trios. There are almost never more than three adult females in a group.<ref>Patry M , Leus K Macdonald AA (1995) Group structure and behaviour of babirusa (Babyrousa babyrussa) in northern Sulawesi. Australian Journal of Zoology 43, 643–655.</ref> The tusks of the adult males are used in intraspecific fighting. The upper tusks are for defense while the lower tusks are offensive weapons. If a male babirusa does not grind his tusks (achievable through regular activity), they can eventually keep growing and, rarely, penetrate the individual's skull.
Female babirusa cycle lengths are between 28 and 42 days and estrus last 2–3 days. The litter size for a babirusa is usually one or two piglets.
Relationship with humans
In Indonesia, the striking appearance of the babirusa has inspired demonic masks.
The Balinese Hindu-era Court of Justice pavilion and the "floating pavilion" of Klungkung palace ruins are notable for painted babirusa raksasa'' (grotesques) on the ceilings.
Prehistoric paintings of babirusa found in caves on the island of Sulawesi in Indonesia have been dated back at least 35,400 years (to the ice age Pleistocene epoch). , who co-authored the 2014 study dating the paintings, said "The paintings of the wild animals are most fascinating because it is clear they were of particular interest to the artists themselves."
The babirusa has sparked debate among Jewish scholars and animal researchers about whether it is considered kosher, or permissible to be consumed by Jews, according to Jewish dietary laws. The debate centers around whether the animal chews its cud, which is a requirement according to the Old Testament for an animal to be considered kosher. Some experts, like J. David Bleich, a professor of Jewish law and ethics at Yeshiva University, believe that the babirusa does not meet the physical criteria to be considered kosher, challenging the assertion that the babirusa chews its cud by citing a report from 1940 that found that true rumination could not take place in the animal's stomach. However, he also notes that Jews can eat any food that is not expressly forbidden and that "the babirusa's resemblance to a pig in appearance and taste is not sufficient grounds for banning its consumption as kosher meat." Others, such as Fuller Bazer, an animal science professor at the University of Florida, believe that the animal is kosher due to its cloven hoof and cud chewing. Additionally, it has been noted that the babirusa is an endangered species and that most Muslims, who face similar dietary restrictions, would avoid eating the meat of any animal whose status in religious law is uncertain.
Conservation status
Babirusas are protected in Indonesia and killing them is illegal in most cases. However, poaching remains a significant threat to the babirusa. Additionally, commercial logging operations threaten the babirusa by habitat loss, and also reduce cover, making the babirusa more exposed to poachers. All extant species of babirusa are listed as vulnerable or endangered by the IUCN.
| Biology and health sciences | Pigs_2 | Animals |
1992713 | https://en.wikipedia.org/wiki/Merlangius | Merlangius | Merlangius merlangus, commonly known as whiting or merling, is an important food fish in the eastern North Atlantic Ocean and the northern Mediterranean, western Baltic, and Black Sea. In Anglophonic countries outside the whiting's natural range, the name "whiting" has been applied to various other species of fish.
Description
Merlangius merlangus has three dorsal fins with a total of 30 to 40 soft rays and two anal fins with 30 to 35 soft rays. The body is long and the head small and a chin barbel, if present, is very small. This fish can reach a maximum length of about . The colour may be yellowish-brown, greenish or dark blue, the flanks yellowish grey or white and the belly silvery. There is a distinctive black blotch near the base of each pectoral fin.
Distribution and habitat
Whiting are native to the northeastern Atlantic Ocean. Their range extends from the southeastern Barents Sea and Iceland to Scandinavia, the Baltic Sea, the North Sea, Portugal, the Black Sea, the Aegean Sea, the Adriatic Sea and parts of the Mediterranean Sea. They occur on sand, mud and gravel seabeds at depths down to about .
In 2014, their conservation status was classified at vulnerable in the Baltic Sea.
Uses
Until the late 20th century, whiting was a cheap fish, regarded as food for the poor or for pets. The general decline in fish stocks means it is now more highly valued. The other fish that have been given the name whiting are mostly also edible fish. Several species of the drum, or croaker, family (Sciaenidae) are also called whiting, among them the northern kingfish (Menticirrhus saxatilis).
Parasites
Whiting and related other Gadidae species are plagued by parasites. These include the cod worm (Lernaeocera branchialis), a copepod crustacean that clings to the gills or the fish and metamorphoses into a plump, sinusoidal, wormlike body, with a coiled mass of egg strings at the rear.
| Biology and health sciences | Acanthomorpha | Animals |
1994190 | https://en.wikipedia.org/wiki/ProPhoto%20RGB%20color%20space | ProPhoto RGB color space | The ProPhoto RGB color space, also known as ROMM RGB (Reference Output Medium Metric), is an output referred RGB color space developed by Kodak. It offers an especially large gamut designed for use with photographic output in mind. The ProPhoto RGB color space encompasses over 90% of possible surface colors in the CIE L*a*b* color space, and 100% of likely occurring real-world surface colors documented by Michael Pointer in 1980, making ProPhoto even larger than the Wide-gamut RGB color space. The ProPhoto RGB primaries were also chosen in order to minimize hue rotations associated with non-linear tone scale operations. One of the downsides to this color space is that approximately 13% of the representable colors are imaginary colors that do not exist and are not visible colors.
When working in color spaces with such a large gamut, it is recommended to work in 16-bit color depth to avoid posterization effects. This will occur more frequently in 8-bit modes as the gradient steps are much larger.
There are two corresponding scene space color encodings known as RIMM RGB (Reference Input Medium Metric) intended to encode standard dynamic range scene space images, and ERIMM RGB intended to encode extended dynamic-range scene space images.
Development
The development of the ProPhoto RGB and other color spaces is documented in an article summarizing a presentation by one of its developers Geoff Wolfe at Kodak, currently senior research manager at Canon Information Systems Research Australia, at the IS&T/SPIE Color Imaging Conference in 2011.
Encoding primaries
Viewing environment
Luminance level is in the range of 160–640 cd/m2.
Viewing surround is average.
There is 0.5–1.0% viewing flare.
The adaptive white point is specified by the chromaticity values for CIE Standard illuminant D50 (, ).
The image color values are assumed to be encoded using flareless (or flare corrected) colorimetric measurements based on the CIE 1931 Standard Colorimetric Observer.
Encoding function
where
and
is the maximum integer value used in the encoding function (e.g. 255 for 8-bit configuration)
and
| Physical sciences | Basics | Physics |
30206738 | https://en.wikipedia.org/wiki/Chronic%20obstructive%20pulmonary%20disease | Chronic obstructive pulmonary disease | Chronic obstructive pulmonary disease (COPD) is a type of progressive lung disease characterized by chronic respiratory symptoms and airflow limitation. GOLD 2024 defined COPD as a heterogeneous lung condition characterized by chronic respiratory symptoms (dyspnea or shortness of breath, cough, sputum production or exacerbations) due to abnormalities of the airways (bronchitis, bronchiolitis) or alveoli (emphysema) that cause persistent, often progressive, airflow obstruction.
The main symptoms of COPD include shortness of breath and a cough, which may or may not produce mucus. COPD progressively worsens, with everyday activities such as walking or dressing becoming difficult. While COPD is incurable, it is preventable and treatable. The two most common types of COPD are emphysema and chronic bronchitis and have been the two classic COPD phenotypes. However, this basic dogma has been challenged as varying degrees of co-existing emphysema, chronic bronchitis, and potentially significant vascular diseases have all been acknowledged in those with COPD, giving rise to the classification of other phenotypes or subtypes.
Emphysema is defined as enlarged airspaces (alveoli) whose walls have broken down resulting in permanent damage to the lung tissue. Chronic bronchitis is defined as a productive cough that is present for at least three months each year for two years. Both of these conditions can exist without airflow limitation when they are not classed as COPD. Emphysema is just one of the structural abnormalities that can limit airflow and can exist without airflow limitation in a significant number of people. Chronic bronchitis does not always result in airflow limitation. However, in young adults with chronic bronchitis who smoke, the risk of developing COPD is high. Many definitions of COPD in the past included emphysema and chronic bronchitis, but these have never been included in GOLD report definitions. Emphysema and chronic bronchitis remain the predominant phenotypes of COPD but there is often overlap between them and a number of other phenotypes have also been described. COPD and asthma may coexist and converge in some individuals. COPD is associated with low-grade systemic inflammation.
The most common cause of COPD is tobacco smoking. Other risk factors include indoor and outdoor air pollution including dust, exposure to occupational irritants such as dust from grains, cadmium dust or fumes, and genetics, such as alpha-1 antitrypsin deficiency. In developing countries, common sources of household air pollution are the use of coal and biomass such as wood and dry dung as fuel for cooking and heating. The diagnosis is based on poor airflow as measured by spirometry.
Most cases of COPD can be prevented by reducing exposure to risk factors such as smoking and indoor and outdoor pollutants. While treatment can slow worsening, there is no conclusive evidence that any medications can change the long-term decline in lung function. COPD treatments include smoking cessation, vaccinations, pulmonary rehabilitation, inhaled bronchodilators and corticosteroids. Some people may benefit from long-term oxygen therapy, lung volume reduction and lung transplantation. In those who have periods of acute worsening, increased use of medications, antibiotics, corticosteroids and hospitalization may be needed.
As of 2015, COPD affected about 174.5 million people (2.4% of the global population). It typically occurs in males and females over the age of 35–40. In 2019 it caused 3.2 million deaths, 80% occurring in lower and middle income countries, up from 2.4 million deaths in 1990. In 2021, it was the fourth biggest cause of death, responsible for approximately 5% of total deaths. The number of deaths is projected to increase further because of continued exposure to risk factors and an aging population. In the United States in 2010 the economic cost was put at US$32.1 billion and projected to rise to US$49 billion in 2020. In the United Kingdom this cost is estimated at £3.8 billion annually.
Signs and symptoms
Shortness of breath
A cardinal symptom of COPD is the chronic and progressive shortness of breath which is most characteristic of the condition. Shortness of breath (breathlessness) is often the most distressing symptom responsible for the associated anxiety and level of disability experienced. Symptoms of wheezing and chest tightness associated with breathlessness can be variable over the course of a day or between days and are not always present. Chest tightness often follows exertion. Many people with more advanced COPD breathe through pursed lips, which can improve shortness of breath. Shortness of breath is often responsible for reduced physical activity and low levels of physical activity are associated with worse outcomes. In severe and very severe cases there may be constant tiredness, weight loss, muscle loss and anorexia. People with COPD often have increased breathlessness and frequent colds before seeking treatment.
Cough
The most often first symptom of COPD is a chronic cough, which may or may not be productive of mucus as phlegm. Phlegm coughed up as sputum can be intermittent and may be swallowed or spat out depending on social or cultural factors and is therefore not always easy to evaluate. However, an accompanying productive cough is only seen in up to 30% of cases. Sometimes limited airflow may develop in the absence of a cough. Symptoms are usually worse in the morning.
A chronic productive cough is the result of mucus hypersecretion and when it persists for more than three months each year for at least two years, it is defined as chronic bronchitis. Chronic bronchitis can occur before the restricted airflow diagnostic of COPD. Some people with COPD attribute the symptoms to the consequences of smoking. In severe COPD, vigorous coughing may lead to rib fractures or to a brief loss of consciousness.
Exacerbations
An acute exacerbation is a sudden worsening of signs and symptoms that lasts for several days. The key symptom is increased breathlessness, other more pronounced symptoms are of excessive mucus, increased cough and wheeze. A commonly found sign is air trapping giving a difficulty in complete exhalation.
The usual cause of an exacerbation is a viral infection, most often the common cold. The common cold is usually associated with the winter months but can occur at any time. Other respiratory infections may be bacterial or in combination sometimes secondary to a viral infection. The most common bacterial infection is caused by Haemophilus influenzae. Other risks include exposure to tobacco smoke (active and passive) and environmental pollutantsboth indoor and outdoor. During the COVID-19 pandemic, hospital admissions for COPD exacerbations sharply decreased which may be attributable to reduction of emissions and cleaner air. There has also been a marked decrease in the number of cold and flu infections during this time.
Smoke from wildfires is proving an increasing risk in many parts of the world and government agencies have published protective advice on their websites. In the US the EPA advises that the use of dust masks do not give protection from the fine particles in wildfires and instead advise the use of well-fitting particulate masks. This same advice is offered in Canada and Australia to the effects of their forest fires.
The number of exacerbations is not seen to relate to any stage of the disease; those with two or more a year are classed as frequent exacerbators and these lead to a worsening in the disease progression. Frailty in ageing increases exacerbations and hospitalization.
Acute exacerbations in COPD are often unexplained and thought to have many causes other than infections. A study has emphasized the possibility of a pulmonary embolism as sometimes being responsible in these cases. Signs can include pleuritic chest pain and heart failure without signs of infection. Such emboli could respond to anticoagulants.
Other conditions
COPD often occurs along with a number of other conditions (comorbidities) due in part to shared risk factors. Common comorbidities include cardiovascular disease, skeletal muscle dysfunction, metabolic syndrome, osteoporosis, depression, anxiety, asthma and lung cancer. Alpha-1 antitrypsin deficiency (A1AD) is an important risk factor for COPD. It is advised that everybody with COPD be screened for A1AD. Metabolic syndrome has been seen to affect up to fifty percent of those with COPD and significantly affects the outcomes. When comorbid with COPD there is more systemic inflammation. It is not known if it co-exists with COPD or develops as a consequence of the pathology. Metabolic syndrome on its own has a high rate of morbidity and mortality and this rate is amplified when comorbid with COPD. Tuberculosis is a risk factor for the development of COPD, and is also a potential comorbidity. Most people with COPD die from comorbidities and not from respiratory problems.
Anxiety and depression are often complications of COPD. Other complications include reduced quality of life and increased disability, cor pulmonale, frequent chest infections including pneumonia, secondary polycythemia, respiratory failure, pneumothorax, lung cancer, and cachexia (muscle wasting).
Along with these complications, there is an associated risk of developing pulmonary hypertension. The estimated prevalence of pulmonary hypertension complicating COPD was reported at 39% in a meta-analysis. Of the people with COPD listed for lung transplantation, 82% were documented as having pulmonary hypertension via right heart catheterization, noting a mean pulmonary arterial pressure greater than 20mm Hg. Despite pulmonary hypertension being relatively rare in people with COPD, mild elevations of pulmonary arterial pressure can lead to worse outcomes, including risk of death.
Cognitive impairment is common in those with COPD as it is for other lung conditions that affect airflow. Cognitive impairment is associated with the declining ability to cope with the basic activities of daily living.
It is unclear if those with COPD are at greater risk of contracting COVID-19, though if infected they are at risk of hospitalization and developing severe COVID-19. However, there are laboratory and clinical studies showing a possibility of certain inhaled corticosteroids for COPD providing a protective role against COVID-19.
Differentiating COVID-19 symptoms from an exacerbation is difficult; mild prodromal symptoms may delay its recognition and where they include loss of taste or smell COVID-19 is to be suspected.
Definition
Many definitions of COPD in the past included chronic bronchitis and emphysema but these have never been included in GOLD report definitions. Emphysema is defined as enlarged airspaces (alveoli) whose walls break down resulting in permanent damage to the lung tissue and is just one of the structural abnormalities that can limit airflow. The condition can exist without airflow limitation but commonly it does. Chronic bronchitis is defined as a productive cough that is present for at least three months each year for two years but does not always result in airflow limitation although the risk of developing COPD is great. These older definitions grouped the two types as type A and type B. Type A were emphysema types known as pink puffers due to their pink complexion, fast breathing rate and pursed lips. Type B were chronic bronchitic types referred to as blue bloaters due to low oxygen levels causing a bluish color to the skin and lips and swollen ankles. These differences were suggested to be due to the presence or not of collateral ventilation, evident in emphysema and lacking in chronic bronchitis. This terminology was no longer accepted as useful, as most people with COPD have a combination of both emphysema and airway disease. These are now recognized as the two major phenotypes of COPD — emphysematous phenotype and chronic bronchitic phenotype.
Subtypes
It has since been recognized that COPD is more complex, with a diverse group of disorders of differing risk factors and clinical courses that has resulted in a number of subtypes or phenotypes of COPD being accepted and proposed. The two classic emphysematous and chronic bronchitic phenotypes are fundamentally different conditions with unique underlying mechanisms. Another subtype of COPD, categorized by some as a separate clinical entity, is asthma-COPD overlap, which is a condition sharing clinical features of both asthma and COPD. Spirometry measures are inadequate for defining phenotypes and chest X-ray, CT and MRI scans have been mostly employed. Most cases of COPD are diagnosed at a late stage and the use of imaging methods would allow earlier detection and treatment.
The identification and recognition of different phenotypes can guide appropriate treatment approaches. For example, the PDE4 inhibitor roflumilast is targeted at the chronic-bronchitic phenotype.
Two inflammatory phenotypes show a phenotype stability: the neutrophilic inflammatory phenotype and the eosinophilic inflammatory phenotype. Mepolizumab, a monoclonal antibody, has been shown to have benefit in treating the eosinophilic inflammatory type rather than the use of oral corticosteroids, but further studies have been called for.
Another recognized phenotype is the frequent exacerbator. The frequent exacerbator has two or more exacerbations a year, has a poor prognosis and is described as a moderately stable phenotype.
A pulmonary vascular COPD phenotype has been described due to cardiovascular dysfunction. A molecular phenotype of CFTR dysfunction is shared with cystic fibrosis. A combined phenotype of chronic bronchitis and bronchiectasis has been described with a difficulty noted of determining the best treatment.
The only genotype is the alpha-1 antitrypsin deficiency (AATD) genetic subtype and this has a specific treatment.
Cause
The most common cause of the development of COPD is the exposure to harmful particles or gases, including tobacco smoke, that irritate the lung causing inflammation that interacts with a number of host factors. Such exposure needs to be significant or long-term. The greatest risk factor for the development of COPD is tobacco smoke. However, less than 50 percent of heavy smokers develop COPD, so other factors need to be considered, including exposure to indoor and outdoor pollutants, allergens, occupational exposure, and host factors. One of the known causes of COPD is exposure to construction dust. The three main types of construction dust are silica dust, non-silica dust (e.g., dust from gypsum, cement, limestone, marble and dolomite) and wood dust. Host factors include a genetic susceptibility, factors associated with poverty, aging and physical inactivity. Asthma and tuberculosis are also recognized as risk factors, as the comorbidity of COPD is reported to be 12 times higher in patients with asthma after adjusting for smoking history. In Europe airway hyperresponsiveness is rated as the second most important risk factor after smoking.
A host factor of an airway branching variation, arising during development has been described. The respiratory tree is a filter for harmful substances and any variant has the potential to disrupt this. A variation has been found to be associated with the development of chronic bronchitis and another with the development of emphysema. A branch variant in the central airway is specifically associated with an increased susceptibility for the later development of COPD. A genetic association for the variants has been sometimes found with FGF10.
Alcohol abuse can lead to alcoholic lung disease and is seen to be an independent risk factor for COPD. Mucociliary clearance is disrupted by chronic exposure to alcohol; macrophage activity is diminished and an inflammatory response promoted. The damage leads to a susceptibility for infection, including COVID-19, more so when combined with smoking; smoking induces the upregulation of the expression of ACE2, a receptor for the SARS-CoV-2 virus.
Smoking
The primary risk factor for COPD globally is tobacco smoking with an increased rate of developing COPD shown in smokers and ex-smokers. Of those who smoke, about 20% will get COPD, increasing to less than 50% in heavy smokers. In the United States and United Kingdom, of those with COPD, 80–95% are either current or previous smokers. Several studies indicate that women are more susceptible than men to the harmful effects of tobacco smoke. For the same amount of cigarette smoking, women have a higher risk of COPD than men. In non-smokers, exposure to second-hand smoke (passive smoking) is the cause of 1.2 million deaths from the more than 8 million deaths worldwide each year due to tobacco smoke. Women who smoke during pregnancy, and during the early life of the child is a risk factor for the later development of COPD in their child.
Inhaled smoke triggers the release of excessive proteases in lungs, which then degrades elastin, the major component of alveoli. Smoke also impairs the action of cilia, inhibiting mucociliary clearance that clears the bronchi of mucus, cellular debris and unwanted fluid.
Other types of tobacco smoke, such as from cigar, pipe, water-pipe and hookah use, also confer a risk. Water-pipe or hookah smoke appears to be as harmful or even more harmful than smoking cigarettes.
Marijuana is the second most commonly smoked substance, but evidence linking its use to COPD is very limited. Limited evidence shows that marijuana does not accelerate lung function decline. A low use of marijuana gives a bronchodilatory effect rather than the bronchoconstrictive effect from tobacco use, but it is often smoked in combination with tobacco or on its own by tobacco smokers. Higher use however has shown a decline in the FEV1. There is evidence of it causing some respiratory problems and its use in combination may have a cumulative toxic effect suggesting it as a risk factor for spontaneous pneumothorax, bullous emphysema, COPD and lung cancer. A noted difference between marijuana use and tobacco was that respiratory problems were resolved with stopping usage unlike the continued decline with stopping tobacco smoking. Respiratory symptoms reported with marijuana use included chronic cough, increased sputum production and wheezing but not shortness of breath. Also these symptoms were typically reported ten years ahead of their affecting tobacco smokers. Another study found that chronic marijuana smokers even with the additional use of tobacco developed similar respiratory problems, but did not seem to develop airflow limitation and COPD.
Pollution
Exposure to particulates can bring about the development of COPD, or its exacerbations. Those with COPD are more susceptible to the harmful effects of particulate exposure that can cause acute exacerbations brought about by infections. Black carbon also known as soot, is an air pollutant associated with an increased risk of hospitalization due to the exacerbations caused. Long-term exposure is indicated as an increased rate of mortality in COPD. Studies have shown that people who live in large cities have a higher rate of COPD compared to people who live in rural areas. Areas with poor outdoor air quality, including that from exhaust gas, generally have higher rates of COPD. Urban air pollution significantly effects the developing lung and its maturation, and contributes a potential risk factor for the later development of COPD. The overall effect in relation to smoking is believed to be small.
Poorly ventilated fires used for cooking and heating, are often fueled by coal or biomass such as wood and dry dung, leading to indoor air pollution and are one of the most common causes of COPD in developing countries. Women are affected more as they have a greater exposure. These fuels are used as the main source of energy in 80% of homes in India, China and sub-Saharan Africa.
Occupational exposure
Intense and prolonged exposure to workplace dusts, chemicals and fumes increases the risk of COPD in smokers, nonsmokers and never-smokers. Substances implicated in occupational exposure and listed in the UK, include organic and inorganic dusts such as cadmium, silica, dust from grains and flour and fumes from cadmium and welding that promote respiratory symptoms. Workplace exposure is believed to be the cause in 10–20% of cases and in the United States, it is believed to be related to around 30% of cases among never smokers and probably represents a greater risk in countries without sufficient regulations. The negative effects of dust exposure and cigarette smoke exposure appear to be cumulative.
Genetics
Genetics play a role in the development of COPD. It is more common among relatives of those with COPD who smoke than unrelated smokers. The most well known genetic risk factor is alpha-1 antitrypsin deficiency (AATD) and this is the only genotype (genetic subtype) with a specific treatment. This risk is particularly high if someone deficient in alpha-1 antitrypsin (AAT) also smokes. It is responsible for about 1–5% of cases and the condition is present in about three to four in 10,000 people.
Mutations in MMP1 gene that encodes for interstitial collagenase are associated with COPD.
The COPDGene study is an ongoing longitudinal study into the epidemiology of COPD, identifying phenotypes and looking for their likely association with susceptible genes. Genome wide analyses in concert with the International COPD Genetics Consortium has identified more than 80 genome regions associated with COPD and further studies in these regions has been called for. Whole genome sequencing is an ongoing
collaboration (2019) with the National Heart, Lung and Blood Institute (NHLBI) to identify rare genetic determinants.
Pathophysiology
COPD is a progressive lung disease in which chronic, incompletely reversible poor airflow (airflow limitation) and an inability to breathe out fully (air trapping) exist. The poor airflow is the result of small airways disease and emphysema (the breakdown of lung tissue). The relative contributions of these two factors vary between people. Air trapping precedes lung hyperinflation.
COPD develops as a significant and chronic inflammatory response to inhaled irritants which ultimately leads to bronchial and alveolar remodelling in the lung known as small airways disease. Thus, airway remodelling with narrowing of peripheral airway and emphysema are responsible for the alteration of lung function. Mucociliary clearance is particularly altered with a dysregulation of cilia and mucus production. Small airway disease sometimes called chronic bronchiolitis, appears to be the precursor for the development of emphysema.
The inflammatory cells involved include neutrophils and macrophages, two types of white blood cells. Those who smoke additionally have cytotoxic T cell involvement and some people with COPD have eosinophil involvement similar to that in asthma. Part of this cell response is brought on by inflammatory mediators such as chemotactic factors. Other processes involved with lung damage include oxidative stress produced by high concentrations of free radicals in tobacco smoke and released by inflammatory cells and breakdown of the connective tissue of the lungs by proteases (particularly elastase) that are insufficiently inhibited by protease inhibitors. The destruction of the connective tissue of the lungs leads to emphysema, which then contributes to the poor airflow and finally, poor absorption and release of respiratory gases. General muscle wasting that often occurs in COPD may be partly due to inflammatory mediators released by the lungs into the blood.
Narrowing of the airways occurs due to inflammation and subsequent scarring within them. This contributes to the inability to breathe out fully. The greatest reduction in air flow occurs when breathing out, as the pressure in the chest is compressing the airways at this time. This can result in more air from the previous breath remaining within the lungs when the next breath is started, resulting in an increase in the total volume of air in the lungs at any given time, a process called air trapping which is closely followed by hyperinflation. Hyperinflation from exercise is linked to shortness of breath in COPD, as breathing in is less comfortable when the lungs are already partly filled. Hyperinflation may also worsen during an exacerbation. There may also be a degree of airway hyperresponsiveness to irritants similar to those found in asthma.
Low oxygen levels and eventually, high carbon dioxide levels in the blood, can occur from poor gas exchange due to decreased ventilation from airway obstruction, hyperinflation and a reduced desire to breathe. During exacerbations, airway inflammation is also increased, resulting in increased hyperinflation, reduced expiratory airflow and worsening of gas transfer. This can lead to low blood oxygen levels which if present for a prolonged period, can result in narrowing of the arteries in the lungs, while emphysema leads to the breakdown of capillaries in the lungs. Both of these conditions may result in pulmonary heart disease also classically known as cor pulmonale.
Diagnosis
The diagnosis of COPD should be considered in anyone over the age of 35 to 40 who has shortness of breath, a chronic cough, sputum production, or frequent winter colds and a history of exposure to risk factors for the disease. Spirometry is then used to confirm the diagnosis.
Spirometry
Spirometry measures the amount of airflow obstruction present and is generally carried out after the use of a bronchodilator, a medication to open up the airways. Two main components are measured to make the diagnosis, the forced expiratory volume in one second (FEV1), which is the greatest volume of air that can be breathed out in the first second of a breath and the forced vital capacity (FVC), which is the greatest volume of air that can be breathed out in a single large breath. Normally, 75–80% of the FVC comes out in the first second and a FEV1/FVC ratio less than 70% in someone with symptoms of COPD defines a person as having the disease. Based on these measurements, spirometry would lead to over-diagnosis of COPD in the elderly. The National Institute for Health and Care Excellence criteria additionally require a FEV1 less than 80% of predicted. People with COPD also exhibit a decrease in diffusing capacity of the lung for carbon monoxide due to decreased surface area in the alveoli, as well as damage to the capillary bed. Testing the peak expiratory flow (the maximum speed of expiration), commonly used in asthma diagnosis, is not sufficient for the diagnosis of COPD.
Screening using spirometry in those without symptoms has uncertain effect and is generally not recommended; however, it is recommended for those without symptoms but with a known risk factor.
Assessment
A number of methods can be used to assess the affects and severity of COPD. The MRC breathlessness scale or the COPD assessment test (CAT) are simple questionnaires that may be used. GOLD refers to a modified MRC scale that if used, needs to include other tests since it is simply a test of breathlessness experienced. Scores on CAT range from 0–40 with the higher the score, the more severe the disease. Spirometry may help to determine the severity of airflow limitation. This is typically based on the FEV1 expressed as a percentage of the predicted "normal" for the person's age, gender, height and weight. Guidelines published in 2011 by American and European medical societies recommend partly basing treatment recommendations on the FEV1. The GOLD guidelines group people into four categories based on symptoms assessment, degree of airflow limitation and history of exacerbations. Weight loss, muscle loss and fatigue are seen in severe and very severe cases.
Use of screening questionnaires, such as COPD diagnostic questionnaire (CDQ), alone or in combination with hand-held flow meters is appropriate for screening of COPD in primary care.
Other tests
A chest X-ray is not useful to establish a diagnosis of COPD but it is of use in either excluding other conditions or including comorbidities such as pulmonary fibrosis and bronchiectasis. Characteristic signs of COPD on X-ray include hyperinflation (shown by a flattened diaphragm and an increased retrosternal air space) and lung hyperlucency. A saber-sheath trachea may also be shown that is indicative of COPD.
A CT scan is not routinely used except for the exclusion of bronchiectasis. An analysis of arterial blood is used to determine the need for oxygen supplementation and assess for high levels of carbon dioxide in the blood; this is recommended in those with an FEV1 less than 35% predicted, those with a peripheral oxygen saturation less than 92% and those with symptoms of congestive heart failure. WHO recommends that all those diagnosed with COPD be screened for alpha-1 antitrypsin deficiency.
Differential diagnosis
COPD may need to be differentiated from other conditions such as congestive heart failure, asthma, bronchiectasis, tuberculosis, obliterative bronchiolitis and diffuse panbronchiolitis. The distinction between asthma and COPD is made on the basis of the symptoms, smoking history and whether airflow limitation is reversible with bronchodilators at spirometry. Chronic bronchitis with normal airflow is not classified as COPD.
Prevention
Most cases of COPD are potentially preventable through decreasing exposure to tobacco smoke and other indoor and outdoor pollutants.
Smoking cessation
The policies of governments, public health agencies and antismoking organizations can reduce smoking rates by discouraging people from starting and encouraging people to stop smoking. Smoking bans in public areas and places of work are important measures to decrease exposure to secondhand smoke and while many places have instituted bans, more are recommended.
In those who smoke, stopping smoking is the only measure shown to slow down the worsening of COPD. Even at a late stage of the disease, it can reduce the rate of worsening lung function and delay the onset of disability and death. Often, several attempts are required before long-term abstinence is achieved. Attempts over 5 years lead to success in nearly 40% of people.
Some smokers can achieve long-term smoking cessation through willpower alone. Smoking, however, is highly addictive and many smokers need further support. The chance of quitting is improved with social support, engagement in a smoking cessation program and the use of medications such as nicotine replacement therapy, bupropion, or varenicline. Combining smoking-cessation medication with behavioral therapy is more than twice as likely to be effective in helping people with COPD stop smoking, compared with behavioral therapy alone.
Occupational health
A number of measures have been taken to reduce the likelihood that workers in at-risk industries—such as coal mining, construction and stonemasonry—will develop COPD. Examples of these measures include the creation of public policy, education of workers and management about the risks, promoting smoking cessation, checking workers for early signs of COPD, use of respirators and dust control. Effective dust control can be achieved by improving ventilation, using water sprays and by using mining techniques that minimize dust generation. If a worker develops COPD, further lung damage can be reduced by avoiding ongoing dust exposure, for example by changing their work role.
Pollution control
Both indoor and outdoor air quality can be improved, which may prevent COPD or slow the worsening of existing disease. This may be achieved by public policy efforts, cultural changes and personal involvement. Many developed countries have successfully improved outdoor air quality through regulations which has resulted in improvements in the lung function of their populations. Individuals are also advised to avoid irritants of indoor and outdoor pollution.
In developing countries one key effort is to reduce exposure to smoke from cooking and heating fuels through improved ventilation of homes and better stoves and chimneys. Proper stoves may improve indoor air quality by 85%. Using alternative energy sources such as solar cooking and electrical heating is also effective. Using fuels such as kerosene or coal might produce less household particulate matter than traditional biomass such as wood or dung, but whether this is better health wise is unclear.
Management
COPD currently has no cure, but the symptoms are treatable and its progression can be delayed, particularly by stopping smoking. The major goals of management are to reduce exposure to risk factors including offering non-pharmacological treatments such as help with stopping smoking. Stopping smoking can reduce the rate of lung function decline and also reduce mortality from smoking-related diseases such as lung cancer and cardiovascular disease. Other recommendations include annual influenza vaccinations and pneumococcal vaccination to help reduce the risk of exacerbations; CDC and GOLD 2024 also recommends RSV vaccine for individuals above 60 years; giving advice as to healthy eating and encouraging physical exercise. Guidance is also advised as to managing breathlessness and stress.
Other illnesses are also managed. An action plan is drawn up and is to be reviewed. Providing people with a personalized action plan, an educational session and support for use of their action plan in the event of an exacerbation, reduces the number of hospital visits and encourages early treatment of exacerbations. When self-management interventions, such as taking corticosteroids and using supplemental oxygen, is combined with action plans, health-related quality of life is improved compared to usual care. In those with COPD who are malnourished, supplementation with vitamin C, vitamin E, zinc and selenium can improve weight, strength of respiratory muscles and health-related quality of life. Significant vitamin D deficiency is common in those with COPD and can cause increased exacerbations. Supplementation when deficient can give a 50% reduction in the number of exacerbations.
A number of medical treatments are used in the management of stable COPD and exacerbations. These include bronchodilators, corticosteroids and antibiotics.
In those with a severe exacerbation, antibiotics improve outcomes. A number of different antibiotics may be used including amoxicillin, doxycycline and azithromycin; whether one is better than the others is unclear. There is no clear evidence of improved outcomes for those with less severe cases. The FDA recommends against the use of fluoroquinolones when other options are available due to higher risks of serious side effects.
In treating acute hypercapnic respiratory failure (acutely raised levels of carbon dioxide), bilevel positive airway pressure (BPAP) can decrease mortality and the need of intensive care. Fewer than 20% of exacerbations require hospital admission. In those without acidosis from respiratory failure, home care may be able to help avoid some admissions.
In those with end-stage disease, palliative care is focused on relieving symptoms. Morphine can improve exercise tolerance. Non-invasive ventilation may be used to support breathing and also reduce daytime breathlessness.
Bronchodilators
Inhaled short-acting bronchodilators are the primary medications used on an as needed basis; their use on a regular basis is not recommended. The two major types are beta2-adrenergic agonists and anticholinergics; either in long-acting or short-acting forms. Beta2–adrenergic agonists target receptors in the smooth muscle cells in bronchioles causing them to relax and allow improved airflow. They reduce shortness of breath, tend to reduce dynamic hyperinflation and improve exercise tolerance. Short-acting bronchodilators have an effect for four hours and for maintenance therapy long acting bronchodilators with an effect of over twelve hours are used. In times of more severe symptoms a short acting agent may be used in combination. An inhaled corticosteroid used with a long-acting beta-2 agonist is more effective than either one on its own.
Which type of long-acting agent, long-acting muscarinic antagonist (LAMA) such as tiotropium or long-acting beta agonist (LABA), is better is unclear and trying each and continuing with the one that works best may be advisable. Both types of agent appear to reduce the risk of acute exacerbations by 15–25%. The combination of LABA/LAMA may reduce COPD exacerbations and improve quality-of-life compared to long-acting bronchodilators alone. The 2018 NICE guideline recommends use of dual long-acting bronchodilators with economic modelling suggesting that this approach is preferable to starting one long acting bronchodilator and adding another later.
Several short-acting β2 agonists are available, including salbutamol (albuterol) and terbutaline. They provide some relief of symptoms for four to six hours. A long-acting beta agonist (LABA) such as salmeterol, formoterol and indacaterol are often used as maintenance therapy. Some feel the evidence of benefits is limited, while others view the evidence of benefit as established. Long-term use appears safe in COPD with adverse effects include shakiness and heart palpitations. When used with inhaled steroids they increase the risk of pneumonia. While steroids and LABAs may work better together, it is unclear if this slight benefit outweighs the increased risks. There is some evidence that combined treatment of LABAs with long-acting muscarinic antagonists (LAMA), an anticholinergic, and LABA +ICS (inhaled corticosteroid) may be similar in benefits in terms of fewer exacerbation's and quality of life measures for moderate to severe COPD, but LAMA+LABA offers better improvements in forced expiratory volume (FEV1%) and a lower risk of pneumonia. All three together, LABA, LAMA and ICS, have some evidence of benefits. Indacaterol requires an inhaled dose once a day and is as effective as the other long-acting β2 agonist drugs that require twice-daily dosing for people with stable COPD.
The two main anticholinergics used in COPD are ipratropium and tiotropium. Ipratropium is a short-acting muscarinic antagonist (SAMA), while tiotropium is long-acting (LAMA). Tiotropium is associated with a decrease in exacerbations and improved quality of life, and tiotropium provides those benefits better than ipratropium. It does not appear to affect mortality or the overall hospitalization rate. Anticholinergics can cause dry mouth and urinary tract symptoms. They are also associated with increased risk of heart disease and stroke. Aclidinium, another long-acting agent, reduces hospitalizations associated with COPD and improves quality of life. The LAMA umeclidinium bromide is another anticholinergic alternative. When compared to tiotropium, the LAMAs aclidinium, glycopyrronium, and umeclidinium appear to have a similar level of efficacy; with all four being more effective than placebo. Further research is needed comparing aclidinium to tiotropium.
Corticosteroids
Inhaled corticosteroids are anti-inflammatories that are recommended by GOLD as a first-line maintenance treatment in COPD cases with repeated exacerbations. Their regular use increases the risk of pneumonia in severe cases. Studies have shown that the risk of pneumonia is associated with all types of corticosteroids; is related to the disease severity and a dose-response relationship has been noted. Oral glucocorticoids can be effective in treating an acute exacerbation. They appear to have fewer side effects than those given intravenously. Five days of steroids work as well as ten or fourteen days.
The use of corticosteroids is associated with a decrease in the number of lymphoid follicles (in the bronchial lymphoid tissue). A triple inhaled therapy of LABA/LAMA/ICS improves lung function, reduces symptoms and exacerbations and is seen to be more effective than mono or dual therapies. NICE guidelines recommend the use of ICSs in people with asthmatic features or features suggesting steroid responsiveness.
PDE4 inhibitors
Phosphodiesterase-4 inhibitors (PDE4 inhibitors) are anti-inflammatories that improve lung function and reduce exacerbations in moderate to severe illness. Roflumilast is a PDE4 inhibitor used orally once daily to reduce inflammation, it has no direct bronchodilatory effects. It is essentially used in treating those with chronic bronchitis along with systemic corticosteroids. Reported adverse effects of roflumilast appear early in treatment, become less with continued treatment and are reversible. One effect is dramatic weight loss and its use is to be avoided in underweight people. It is also advised to be used with caution in those who have depression.
Other medications
Long-term preventive use of antibiotics, specifically those from the macrolide class such as erythromycin, reduce the frequency of exacerbations in those who have two or more a year. This practice may be cost effective in some areas of the world. Concerns include the potential for antibiotic resistance and side effects including hearing loss, tinnitus and changes to the heart rhythm known as long QT syndrome.
Methylxanthines such as theophylline are widely used. Theophylline is seen to have a mild bronchodilatory effect in stable COPD. Inspiratory muscle function is seen to be improved but the causal effect is unclear. Theophylline is seen to improve breathlessness when used as an add-on to salmeterol. All instances of improvement have been reported using sustained release preparations. Methylxanthines are not recommended for use in exacerbations due to adverse effects.
Mucolytics may help to reduce exacerbations in some people with chronic bronchitis; noticed by less hospitalization and less days of disability in one month. Erdosteine is recommended by NICE. GOLD also supports the use of some mucolytics that are advised against when inhaled corticosteroids are being used and singles out erdosteine as having good effects regardless of corticosteroid use. Erdosteine also has antioxidant properties but there is not enough evidence to support the general use of antioxidants. Erdosteine has been shown to significantly reduce the risk of exacerbations, shorten their duration and hospital stays.
Cough medicines are not recommended. Beta blockers are not contraindicated for those with COPD and should only be used where there is concomitant cardiovascular disease.
Recent studies show that metformin plays a role in reducing systemic inflammation by reducing biomarker levels that are increased during COPD exacerbations.
Oxygen therapy
Supplemental oxygen is recommended for those with low oxygen levels in respiratory failure at rest (a partial pressure of oxygen less than 50–55 mmHg or oxygen saturations of less than 88%). When taking into account complications including cor pulmonale and pulmonary hypertension, the levels involved are 56–59 mmHg. Oxygen therapy is to be used for between 15 and 18 hours per day and is said to decrease the risk of heart failure and death. In those with normal or mildly low oxygen levels, oxygen supplementation (ambulatory) may improve shortness of breath when given during exercise, but may not improve breathlessness during normal daily activities or affect the quality of life. During acute exacerbations, many require oxygen therapy; the use of high concentrations of oxygen without taking into account a person's oxygen saturations may lead to increased levels of carbon dioxide and worsened outcomes. In those at high risk of high carbon dioxide levels, oxygen saturations of 88–92% are recommended, while for those without this risk, recommended levels are 94–98%. Once prescribed long-term oxygen therapy, patients should be re-assessed after 60 to 90 days, to deteermine whether supplemental oxygen is still indicated and if prescribed supplemental oxygen is effective.
Rehabilitation
Pulmonary rehabilitation is a program of exercise, disease management and counseling, coordinated to benefit the individual. A severe exacerbation leads to hospital admission, high mortality and a decline in the ability to carry out daily activities. Following a hospital admission pulmonary rehabilitation has been shown to significantly reduce future hospital admissions, mortality and improve quality of life.
The optimal exercise routine, use of noninvasive ventilation during exercise and intensity of exercise suggested for people with COPD, is unknown. Performing endurance arm exercises improves arm movement for people with COPD and may result in a small improvement in breathlessness. Performing arm exercises alone does not appear to improve quality of life. Pursed-lip breathing exercises may be useful. Tai chi exercises appear to be safe to practice for people with COPD and may be beneficial for pulmonary function and pulmonary capacity when compared to a regular treatment program. Tai Chi was not found to be more effective than other exercise intervention programs. Inspiratory and expiratory muscle training (IMT, EMT) have been suggested and may provide some improvements when compared to no treatment. A combination of IMT and walking exercises at home may help limit breathlessness in cases of severe COPD. Additionally, the use of low amplitude high velocity joint mobilization together with exercise improves lung function and exercise capacity. The goal of spinal manipulation therapy is to improve thoracic mobility in an effort to reduce the work on the lungs during respiration, however, the evidence supporting manual therapy for people with COPD is very weak.
Airway clearance techniques (ACTs), such as postural drainage, percussion/vibration, autogenic drainage, hand-held positive expiratory pressure (PEP) devices and other mechanical devices, may reduce the need for increased ventilatory assistance, the duration of ventilatory assistance and the length of hospital stay in people with acute COPD. In people with stable COPD, ACTs may lead to short-term improvements in health-related quality of life and a reduced long-term need for hospitalizations related to respiratory issues.
Being either underweight or overweight can affect the symptoms, degree of disability and prognosis of COPD. People with COPD who are underweight can improve their breathing muscle strength by increasing their calorie intake. When combined with regular exercise or a pulmonary rehabilitation program, this can lead to improvements in COPD symptoms. Supplemental nutrition may be useful in those who are malnourished.
Management of exacerbations
People with COPD can experience exacerbations (flare-ups) that are commonly caused by respiratory tract infections. The symptoms that worsen are not specific to COPD and differential diagnoses need to be considered. Acute exacerbations are typically treated by increasing the use of short-acting bronchodilators including a combination of a short-acting inhaled beta agonist and short-acting anticholinergic. These medications can be given either via a metered-dose inhaler with a spacer or via a nebulizer, with both appearing to be equally effective. Nebulization may be easier for those who are more unwell. Oxygen supplementation can be useful. Excessive oxygen; however, can result in increased levels and a decreased level of consciousness. Corticosteroids given orally can improve lung function and shorten hospital stays but their use is recommended for only five to seven days; longer courses increase the risk of pneumonia and death.
Room temperature
Maintaining room temperature of at least for a minimum of nine hours a day was associated with better health in those with COPD, especially for smokers. The World Health Organization (WHO) recommends indoor temperatures of a slightly higher range between .
Room humidity
For people with COPD, the ideal indoor humidity levels are 30–50% RH. Maintaining indoor humidity can be difficult in the winter, especially in cold climates where the heating system is constantly running.
Keeping the indoor relative humidity above 40% RH significantly reduces the infectivity of aerosolized viruses.
Procedures for emphysema
There are a number of procedures to reduce the volume of a lung in cases of severe emphysema with hyperinflation.
Surgical
For severe emphysema that has proved unresponsive to other therapies lung volume reduction surgery (LVRS) may be an option. LVRS involves the removal of damaged tissue, which improves lung function by allowing the rest of the lungs to expand. It is considered when the emphysema is in the upper lobes and when there are no comorbidities.
Bronchoscopic
Minimally invasive bronchoscopic procedures may be carried out to reduce lung volume. These include the use of valves, coils, or thermal ablation. Endobronchial valves are one-way valves that may be used in those with severe hyperinflation resulting from advanced emphysema; a suitable target lobe and no collateral ventilation are required for this procedure. The placement of one or more valves in the lobe induces a partial collapse of the lobe that ensures a reduction in residual volume that improves lung function, the capacity for exercise and quality of life.
The placement of nitinol coils instead of valves is recommended where there is collateral ventilation that would prevent the use of valves. Nitinol is a biocompatible alloy.
Both of these techniques are associated with adverse effects including persistent air leaks and cardiovascular complications. Thermal vapor ablation has an improved profile. Heated water vapor is used to target lobe regions which leads to permanent fibrosis and volume reduction. The procedure is able to target individual lobe segments, can be carried out regardless of collateral ventilation and can be repeated with the natural advance of emphysema.
Other surgeries
In very severe cases lung transplantation might be considered. A CT scan may be useful in surgery considerations. Ventilation/perfusion scintigraphy is another imaging method that may be used to evaluate cases for surgical interventions and also to evaluate post-surgery responses. A bullectomy may be carried out when a giant bulla occupies more than a third of a hemithorax.
Prognosis
COPD is progressive and can lead to premature death. It is estimated that 3% of all disability is related to COPD. The proportion of disability from COPD globally has decreased from 1990 to 2010 due to improved indoor air quality primarily in Asia. The overall number of years lived with disability from COPD, however, has increased.
There are many variables affecting the long-term outcome in COPD and GOLD recommends the use of a composite test (BODE) that includes the main variables of body-mass index, obstruction of airways, dyspnea (breathlessness) and exercise and not just spirometry results.
NICE recommends against the use of BODE for the prognosis assessment in stable COPD; factors such as exacerbations and frailty need to be considered. Other factors that contribute to a poor outcome include older age, comorbidities such as lung cancer and cardiovascular disease and the number and severity of exacerbations needing hospital admittance.
Epidemiology
Estimates of prevalence have considerable variation due to differences in analytical and surveying approach and the choice of diagnostic criteria. An estimated 384 million people aged 30 years or more had COPD in 2010, corresponding to a global prevalence of 12%. The disease affects men and women. The increase in the developing world between 1970 and the 2000s is believed to be related to increasing rates of smoking in this region, an increasing population and an aging population due to fewer deaths from other causes such as infectious diseases. Some developed countries have seen increased rates, some have remained stable and some have seen a decrease in COPD prevalence.
Around three million people die of COPD each year. In some countries, mortality has decreased in men but increased in women. This is most likely due to rates of smoking in women and men becoming more similar. A higher rate of COPD is found in those over 40 years and this increases greatly with advancing age with the highest rate found in those over 60 years. Sex differences in the anatomy of the respiratory system include smaller airway lumens and thicker airway walls in women, which contribute to a greater severity of COPD symptoms like dyspnea and frequency of COPD exacerbation.
In the UK, three million people are reported to be affected by COPDtwo million of these being undiagnosed. On average, the number of COPD-related deaths between 2007 and 2016 was 28,600. The estimated number of deaths due to occupational exposure was estimated to be about 15% at around 4,000. In the United States in 2018, almost 15.7 million people had been diagnosed with COPD and it is estimated that millions more have not been diagnosed.
In 2011, there were approximately 730,000 hospitalizations in the United States for COPD. Globally, COPD in 2019 was the third-leading cause of death. In low-income countries, COPD does not appear in the Top 10 causes of death; in other income groups, it is in the Top 5.
History
The name chronic obstructive pulmonary disease is believed to have first been used in 1965. Previously it has been known by a number of different names, including chronic obstructive bronchopulmonary disease, chronic airflow obstruction, chronic obstructive lung disease, nonspecific chronic pulmonary disease, diffuse obstructive pulmonary syndrome.
The terms emphysema and chronic bronchitis were formally defined as components of COPD in 1959 at the CIBA guest symposium and in 1962 at the American Thoracic Society Committee meeting on Diagnostic Standards.
Early descriptions of probable emphysema began in 1679 by T. Bonet of a condition of "voluminous lungs" and in 1769 by Giovanni Morgagni of lungs which were "turgid particularly from air". In 1721 the first drawings of emphysema were made by Ruysh. René Laennec, used the term emphysema in his book A Treatise on the Diseases of the Chest and of Mediate Auscultation (1837) to describe lungs that did not collapse when he opened the chest during an autopsy. He noted that they did not collapse as usual because they were full of air and the airways were filled with mucus. In 1842, John Hutchinson invented the spirometer, which allowed the measurement of vital capacity of the lungs. However, his spirometer could only measure volume, not airflow. Tiffeneau and Pinelli in 1947 described the principles of measuring airflow.
Air pollution and the increase in cigarette smoking in Great Britain at the start of the 20th century led to high rates of chronic lung disease, though it received little attention until the Great Smog of London in December 1952. This spurred epidemiological research in the United Kingdom, Holland and elsewhere. In 1953, George L. Waldbott, an American allergist, first described a new disease he named smoker's respiratory syndrome in the 1953 Journal of the American Medical Association. This was the first association between tobacco smoking and chronic respiratory disease.
Modern treatments were developed during the second half of the 20th century. Evidence supporting the use of steroids in COPD was published in the late 1950s. Bronchodilators came into use in the 1960s following a promising trial of isoprenaline. Further bronchodilators, such as short-acting salbutamol, were developed in the 1970s and the use of long-acting bronchodilators began in the mid-1990s.
Society and culture
It is generally accepted that COPD is widely underdiagnosed and many people remain untreated. In the US the NIH has promoted November as COPD Awareness Month to be an annual focus on increasing awareness of the condition.
Economics
Globally, as of 2010, COPD is estimated to result in economic costs of $2.1 trillion, half of which occurring in the developing world. Of this total an estimated $1.9 trillion are direct costs such as medical care, while $0.2 trillion are indirect costs such as missed work. This is expected to more than double by 2030. In Europe, COPD represents 3% of healthcare spending. In the United States, costs of the disease are estimated at $50 billion, most of which is due to exacerbation. COPD was among the most expensive conditions seen in U.S. hospitals in 2011, with a total cost of about $5.7 billion.
Research
Hyaluronan is a natural sugar in the extracellular matrix that provides a protective coating for cells. It has been shown that on exposure to pollution the hyaluronan in the lungs breaks down into fragments causing irritation and activation of the immune system. There follows subsequent airway constriction and inflammation. The study showed that the inhalation of unfragmented hyaluronan overcame the effects of fragmented HA and reduced inflammation. Inhaled HA only acts locally in the bronchial tree and does not interfere with any drug. It improves mucus clearance by allowing it to move more freely. Further studies are to be carried out in the
US to determine optimum dosage levels.
A new cryogenic treatment aimed at the chronic bronchitic subtype using a liquid nitrogen metered cryospray is being trialed and was due to complete in September 2021.
Stem-cell therapy using mesenchymal stem cells has the potential to restore lung function and thereby improve quality of life. In June 2021 eight clinical trials had been completed and seventeen were underway. Overall, stem cell therapy has proved to be safe. The trials include the use of stem cells from different sources such as adipose tissue, bone marrow and umbilical cord blood.
A procedure known as targeted lung denervation is being trialed and has been used as part of a clinical trial (2021) in a hospital in the UK. The new minimally invasive procedure which takes about an hour to carry out, places electrodes to destroy branches of the vagus nerve in the lungs. The vagus nerve is responsible for both muscle contraction and mucus secretion, which results in narrowing the airways. In those with COPD these nerves are overactive, usually as a result of smoking damage and the constant mucus secretion and airway constriction leads to the symptoms of cough, shortness of breath, wheeze and tightness of the chest.
The effectiveness of alpha-1 antitrypsin augmentation treatment for people who have alpha-1 antitrypsin deficiency is unclear. A later clinical trial of double-dosing has shown some improvements in slowing the breakdown of elastin and the progression of emphysema with further studies being called for.
Mass spectrometry is being studied as a diagnostic tool in COPD.
Research continues into the use of telehealthcare to treat people with COPD when they experience episodes of shortness of breath; treating people remotely may reduce the number of emergency-room visits and improve the person's quality of life.
Evidence is growing for the effectiveness of Astaxanthin against lung disease including COPD. Astaxanthin is a potent antioxidant with anti-inflammatory properties and more trials are said to be needed into its use.
American COPD patients and their caregivers consider the following COPD-related research areas as the most important:
family/social/community research
patients' well-being
curative research
biomedical therapies
policy
holistic therapies.
Other animals
Chronic obstructive pulmonary disease may occur in a number of other animals and may be caused by exposure to tobacco smoke. Most cases of the disease, however, are relatively mild. In horses it is known as recurrent airway obstruction (RAO) or heaves. RAO can be quite severe and most often is linked to exposure to common allergens. COPD is also commonly found in old dogs.
| Biology and health sciences | Non-infectious disease | null |
1354990 | https://en.wikipedia.org/wiki/NGC%201275 | NGC 1275 | NGC 1275 (also known as Perseus A or Caldwell 24) is a type 1.5 Seyfert galaxy located around 237 million light-years away in the direction of the constellation Perseus. NGC 1275 is a member of the large Perseus Cluster of galaxies.
Properties
NGC 1275 consists of two galaxies, a central type-cD galaxy in the Perseus Cluster, and a so-called high velocity system (HVS) which lies in front of it. The HVS is moving at 3000 km/s towards the dominant system, and is believed to be merging with the Perseus Cluster. The HVS is not affecting the cD galaxy as it lies at least 200 thousand light years from it. However tidal interactions are disrupting it and the ram pressure produced by its interaction with the intracluster medium of Perseus is stripping its gas as well as producing large amounts of star formation within it.
The central cluster galaxy contains a massive network of spectral line emitting filaments, which apparently are being dragged out by rising bubbles of relativistic plasma generated by the central active galactic nucleus. Long gaseous filaments made up of threads of gas stretch out beyond the galaxy, into the multimillion-degree, X-ray–emitting gas that fills the cluster. The amount of gas contained in a typical thread is approximately one million times the mass of the Sun. They are only 200 light-years wide, are often very straight, and extend for up to 20,000 light-years.
The existence of the filaments poses a problem. As they are much cooler than the surrounding intergalactic cloud, it is unclear how they have existed for such a long time, or why they have not warmed, dissipated or collapsed to form stars. One possibility is that weak magnetic fields (about one-ten-thousandth the strength of Earth's field) exert enough force on the ions within the threads to keep them together.
NGC 1275 contains 13 billion solar masses of molecular hydrogen that seems to be infalling from Perseus' intracluster medium in a cooling flow, both feeding its active nucleus and fueling significant amounts of star formation
The presence of an active nucleus demonstrates that a supermassive black hole is present in NGC 1275's center. The black hole is surrounded by a rotating disk of molecular gas. High-resolution observations of the rotation of this disk obtained using adaptive optics at the Gemini North telescope indicate a central mass of approximately 800 million solar masses, including both the mass of the black hole and of the inner core of the gas disk.
Supernovae
Three supernovae have been observed in NGC 1275:
SN 1968A (type unknown, mag. 15.5) was discovered by Miklos Lovas on 25 January 1968.
SN 2005mz (type Ia, mag. 18.2) was discovered by Jack Newton, M. Peoples, and Tim Puckett on 31 December 2005.
SN 2024xav (typeII-P, mag. 18.63) was discovered by GOTO on 2 October 2024.
| Physical sciences | Notable galaxies | Astronomy |
1355057 | https://en.wikipedia.org/wiki/Overhead%20power%20line | Overhead power line | An overhead power line is a structure used in electric power transmission and distribution to transmit electrical energy along large distances. It consists of one or more conductors (commonly multiples of three) suspended by towers or poles. Since the surrounding air provides good cooling, insulation along long passages, and allows optical inspection, overhead power lines are generally the lowest-cost method of power transmission for large quantities of electric energy.
Construction
Towers for support of the lines are made of wood (as-grown or laminated), steel or aluminum (either lattice structures or tubular poles), concrete, and occasionally reinforced plastics. The bare wire conductors on the line are generally made of aluminum (either plain or reinforced with steel, or composite materials such as carbon and glass fiber), though some copper wires are used in medium-voltage distribution and low-voltage connections to customer premises. A major goal of overhead power line design is to maintain adequate clearance between energized conductors and the ground so as to prevent dangerous contact with the line, and to provide reliable support for the conductors, resilience to storms, ice loads, earthquakes and other potential damage causes.
Today some overhead lines are routinely operated at voltages exceeding 765,000 volts between conductors, with even higher voltages possible in some cases.
Classification by operating voltage
Overhead power transmission lines are classified in the electrical power industry by the range of voltages:
Low voltage (LV) – less than 1000 Volts, used for connection between a residential or small commercial customer and the utility.
Medium voltage (MV; distribution) – between 1000 Volts (1 kV) and 69 kV, used for distribution in urban and rural areas.
High voltage (HV; subtransmission less than 100 kV; subtransmission or transmission at voltages such as 115 kV and 138 kV), used for sub-transmission and transmission of bulk quantities of electric power and connection to very large consumers.
Extra high voltage (EHV; transmission) – from 345 kV, up to about 800 kV, used for long distance, very high power transmission.
Ultra high voltage (UHV) – higher than 800 kV. The Financial Times reported UHV lines are a "game changer", making a global electricity grid potentially feasible. StateGrid said that compared to conventional lines, UHV enables the transmission of five times more power, over six times the distance.
Structures
Structures for overhead lines take a variety of shapes depending on the type of line. Structures may be as simple as wood poles directly set in the earth, carrying one or more cross-arm beams to support conductors, or "armless" construction with conductors supported on insulators attached to the side of the pole. Tubular steel poles are typically used in urban areas. High-voltage lines are often carried on lattice-type steel towers or pylons. For remote areas, aluminum towers may be placed by helicopters. Concrete poles have also been used. Poles made of reinforced plastics are also available, but their high cost restricts application.
Each structure must be designed for the loads imposed on it by the conductors. The weight of the conductor must be supported, as well as dynamic loads due to wind and ice accumulation, and effects of vibration. Where conductors are in a straight line, towers need only resist the weight since the tension in the conductors approximately balances with no resultant force on the structure. Flexible conductors supported at their ends approximate the form of a catenary, and much of the analysis for construction of transmission lines relies on the properties of this form.
A large transmission line project may have several types of towers, with "tangent" ("suspension" or "line" towers, UK) towers intended for most positions and more heavily constructed towers used for turning the line through an angle, dead-ending (terminating) a line, or for important river or road crossings. Depending on the design criteria for a particular line, semi-flexible type structures may rely on the weight of the conductors to be balanced on both sides of each tower. More rigid structures may be intended to remain standing even if one or more conductors is broken. Such structures may be installed at intervals in power lines to limit the scale of cascading tower failures.
Foundations for tower structures may be large and costly, particularly if the ground conditions are poor, such as in wetlands. Each structure may be stabilized considerably by the use of guy wires to counteract some of the forces applied by the conductors.
Power lines and supporting structures can be a form of visual pollution. In some cases the lines are buried to avoid this, but this "undergrounding" is more expensive and therefore not common.
For a single wood utility pole structure, a pole is placed in the ground, then three crossarms extend from this, either staggered or all to one side. The insulators are attached to the crossarms. For an "H"-type wood pole structure, two poles are placed in the ground, then a crossbar is placed on top of these, extending to both sides. The insulators are attached at the ends and in the middle. Lattice tower structures have two common forms. One has a pyramidal base, then a vertical section, where three crossarms extend out, typically staggered. The strain insulators are attached to the crossarms. Another has a pyramidal base, which extends to four support points. On top of this a horizontal truss-like structure is placed.
A grounded wire is sometimes strung along the tops of the towers to provide lightning protection. An optical ground wire is a more advanced version with embedded optical fibers for communication. Overhead wire markers can be mounted on the ground wire to meet International Civil Aviation Organization recommendations.
Some markers include flashing lamps for night-time warning.
Circuits
A single-circuit transmission line carries conductors for only one circuit. For a three-phase system, this implies that each tower supports three conductors.
A double-circuit transmission line has two circuits. For three-phase systems, each tower supports and insulates six conductors. Single phase AC-power lines as used for traction current have four conductors for two circuits. Usually both circuits operate at the same voltage.
In HVDC systems typically two conductors are carried per line, but in rare cases only one pole of the system is carried on a set of towers.
In some countries like Germany most power lines with voltages above 100 kV are implemented as double, quadruple or in rare cases even hextuple power line as rights of way are rare. Sometimes all conductors are installed with the erection of the pylons; often some circuits are installed later. A disadvantage of double circuit transmission lines is that maintenance can be difficult, as either work in close proximity of high voltage or switch-off of two circuits is required. In case of failure, both systems can be affected.
The largest double-circuit transmission line is the Kita-Iwaki Powerline.
Insulators
Insulators must support the conductors and withstand both the normal operating voltage and surges due to switching and lightning. Insulators are broadly classified as either pin-type, which support the conductor above the structure, or suspension type, where the conductor hangs below the structure. The invention of the strain insulator was a critical factor in allowing higher voltages to be used.
At the end of the 19th century, the limited electrical strength of telegraph-style pin insulators limited the voltage to no more than 69,000 volts. Up to about 33 kV (69 kV in North America) both types are commonly used. At higher voltages only suspension-type insulators are common for overhead conductors.
Insulators are usually made of wet-process porcelain or toughened glass, with increasing use of glass-reinforced polymer insulators. However, with rising voltage levels, polymer insulators (silicone rubber based) are seeing increasing usage. China has already developed polymer insulators having a highest system voltage of 1100 kV and India is currently developing a 1200 kV (highest system voltage) line which will initially be charged with 400 kV to be upgraded to a 1200 kV line.
Suspension insulators are made of multiple units, with the number of unit insulator disks increasing at higher voltages. The number of disks is chosen based on line voltage, lightning withstand requirement, altitude, and environmental factors such as fog, pollution, or salt spray. In cases where these conditions are suboptimal, longer insulators must be used. Longer insulators with longer creepage distance for leakage current, are required in these cases. Strain insulators must be strong enough mechanically to support the full weight of the span of conductor, as well as loads due to ice accumulation, and wind.
Porcelain insulators may have a semi-conductive glaze finish, so that a small current (a few milliamperes) passes through the insulator. This warms the surface slightly and reduces the effect of fog and dirt accumulation. The semiconducting glaze also ensures a more even distribution of voltage along the length of the chain of insulator units.
Polymer insulators by nature have hydrophobic characteristics providing for improved wet performance. Also, studies have shown that the specific creepage distance required in polymer insulators is much lower than that required in porcelain or glass. Additionally, the mass of polymer insulators (especially in higher voltages) is approximately 50% to 30% less than that of a comparative porcelain or glass string. Better pollution and wet performance is leading to the increased use of such insulators.
Insulators for very high voltages, exceeding 200 kV, may have grading rings installed at their terminals. This improves the electric field distribution around the insulator and makes it more resistant to flash-over during voltage surges.
Conductors
The most common conductor in use for transmission today is aluminum conductor steel reinforced (ACSR). Also seeing much use is all-aluminum-alloy conductor (AAAC). Aluminum is used because it has about half the weight of a comparable resistance copper cable (though larger diameter due to lower specific conductivity), as well as being cheaper.
Copper was more popular in the past and is still in use, especially at lower voltages and for grounding.
While larger conductors lose less energy due to lower electrical resistance, they are more costly than smaller conductors. An optimization rule called Kelvin's Law (named for Lord Kelvin) states that the optimum size of conductor for a line is found when the cost of the energy wasted in the conductor is equal to the annual interest paid on that portion of the line construction cost due to the size of the conductors. The optimization problem is made more complex by additional factors such as varying annual load, varying cost of installation, and the discrete sizes of cable that are commonly made.
Since a conductor is a flexible object with uniform weight per unit length, the shape of a conductor strung between two towers approximates that of a catenary. The sag of the conductor (vertical distance between the highest and lowest point of the curve) varies depending on the temperature and additional load such as ice cover. A minimum overhead clearance must be maintained for safety. Since the temperature and therefore length of the conductor increase with increasing current through it, it is sometimes possible to increase the power handling capacity (uprate) by changing the conductors for a type with a lower coefficient of thermal expansion or a higher allowable operating temperature.
Two such conductors that offer reduced thermal sag are known as composite core conductors (ACCR and ACCC conductor). In lieu of steel core strands that are often used to increase overall conductor strength, the ACCC conductor uses a carbon and glass fiber core that offers a coefficient of thermal expansion about 1/10 of that of steel. While the composite core is nonconductive, it is substantially lighter and stronger than steel, which allows the incorporation of 28% more aluminum (using compact trapezoidal-shaped strands) without any diameter or weight penalty. The added aluminum content helps reduce line losses by 25 to 40% compared to other conductors of the same diameter and weight, depending upon electric current. The carbon core conductor's reduced thermal sag allows it to carry up to twice the current ("ampacity") compared to all-aluminum conductor (AAC) or ACSR.
The power lines and their surroundings must be maintained by linemen, sometimes assisted by helicopters with pressure washers or circular saws which may work three times faster. However this work often occurs in the dangerous areas of the Helicopter height–velocity diagram, and the pilot must be qualified for this "human external cargo" method.
Bundle conductors
For transmission of power across long distances, high voltage transmission is employed. Transmission higher than 132 kV poses the problem of corona discharge, which causes significant power loss and interference with communication circuits. To reduce this corona effect, it is preferable to use more than one conductor per phase, or bundled conductors.
Bundle conductors consist of several parallel cables connected at intervals by spacers, often in a cylindrical configuration. The optimum number of conductors depends on the current rating, but typically higher-voltage lines also have higher current. American Electric Power is building 765 kV lines using six conductors per phase in a bundle. Spacers must resist the forces due to wind, and magnetic forces during a short-circuit.
Bundled conductors reduce the voltage gradient in the vicinity of the line. This reduces the possibility of corona discharge. At extra high voltage, the electric field gradient at the surface of a single conductor is high enough to ionize air, which wastes power, generates unwanted audible noise and interferes with communication systems. The field surrounding a bundle of conductors is similar to the field that would surround a single, very large conductor—this produces lower gradients which mitigates issues associated with high field strength. The transmission efficiency is improved as loss due to corona effect is countered.
Bundled conductors cool themselves more efficiently due to the increased surface area of the conductors, further reducing line losses. When transmitting alternating current, bundle conductors also avoid the reduction in ampacity of a single large conductor due to the skin effect. A bundle conductor also has lower reactance, compared to a single conductor.
While wind resistance is higher, wind-induced oscillation can be damped at bundle spacers. The ice and wind loading of bundled conductors will be greater than a single conductor of the same total cross section, and bundled conductors are more difficult to install than single conductors.
Ground wires
Overhead power lines are often equipped with a ground conductor (shield wire, static wire, or overhead earth wire). The ground conductor is usually grounded (earthed) at the top of the supporting structure, to minimize the likelihood of direct lightning strikes to the phase conductors. In circuits with earthed neutral, it also serves as a parallel path with the earth for fault currents. Very high-voltage transmission lines may have two ground conductors. These are either at the outermost ends of the highest cross beam, at two V-shaped mast points, or at a separate cross arm. Older lines may use surge arresters every few spans in place of a shield wire; this configuration is typically found in the more rural areas of the United States. By protecting the line from lightning, the design of apparatus in substations is simplified due to lower stress on insulation. Shield wires on transmission lines may include optical fibers (optical ground wires/OPGW), used for communication and control of the power system.
At some HVDC converter stations, the ground wire is used also as the electrode line to connect to a distant grounding electrode. This allows the HVDC system to use the earth as one conductor. The ground conductor is mounted on small insulators bridged by lightning arrestors above the phase conductors. The insulation prevents electrochemical corrosion of the pylon.
Medium-voltage distribution lines may also use one or two shield wires, or may have the grounded conductor strung below the phase conductors to provide some measure of protection against tall vehicles or equipment touching the energized line, as well as to provide a neutral line in Wye wired systems.
On some power lines for very high voltages in the former Soviet Union, the ground wire is used for PLC systems and mounted on insulators at the pylons.
Insulated conductors and cable
Overhead insulated cables are rarely used, usually for short distances (less than a kilometer). Insulated cables can be directly fastened to structures without insulating supports. An overhead line with bare conductors insulated by air is typically less costly than a cable with insulated conductors.
A more common approach is "covered" line wire. It is treated as bare cable, but often is safer for wildlife, as the insulation on the cables increases the likelihood of a large-wing-span raptor to survive a brush with the lines, and reduces the overall danger of the lines slightly. These types of lines are often seen in the eastern United States and in heavily wooded areas, where tree-line contact is likely. The only pitfall is cost, as insulated wire is often costlier than its bare counterpart. Many utility companies implement covered line wire as jumper material where the wires are often closer to each other on the pole, such as an underground riser/pothead, and on reclosers, cutouts and the like.
Dampers
Because power lines can suffer from aeroelastic flutter driven by wind, Stockbridge dampers are often attached to the lines to reduce the vibrations.
Compact transmission lines
A compact overhead transmission line requires a smaller right of way than a standard overhead powerline. Conductors must not get too close to each other. This can be achieved either by short span lengths and insulating crossbars, or by separating the conductors in the span with insulators. The first type is easier to build as it does not require insulators in the span, which may be difficult to install and to maintain.
Examples of compact lines are:
Lutsk compact overhead powerline
Hilpertsau-Weisenbach compact overhead line
Compact transmission lines may be designed for voltage upgrade of existing lines to increase the power that can be transmitted on an existing right of way.
Low voltage
Low voltage overhead lines may use either bare conductors carried on glass or ceramic insulators or an aerial bundled cable system. The number of conductors may be anywhere between two (most likely a phase and neutral) up to as many as six (three phase conductors, separate neutral and earth plus street lighting supplied by a common switch); a common case is four (three phase and neutral, where the neutral might also serve as a protective earthing conductor).
Train power
Overhead lines or overhead wires are used to transmit electrical energy to trams, trolleybuses or trains. Overhead line is designed on the principle of one or more overhead wires situated over rail tracks. Feeder stations at regular intervals along the overhead line supply power from the high-voltage grid. For some cases low-frequency AC is used, and distributed by a special traction current network.
Further applications
Overhead lines are also occasionally used to supply transmitting antennas, especially for efficient transmission of long, medium and short waves. For this purpose a staggered array line is often used. Along a staggered array line the conductor cables for the supply of the earth net of the transmitting antenna are attached on the exterior of a ring, while the conductor inside the ring, is fastened to insulators leading to the high-voltage standing feeder of the antenna.
Use of area under overhead power lines
Use of the area below an overhead line is limited because objects must not come too close to the energized conductors. Overhead lines and structures may shed ice, creating a hazard. Radio reception can be impaired under a power line, due both to shielding of a receiver antenna by the overhead conductors, and by partial discharge at insulators and sharp points of the conductors which creates radio noise.
In the area surrounding the overhead lines it is dangerous to risk interference; e.g. flying kites or balloons, using ladders or operating machinery.
Overhead distribution and transmission lines near airfields are often marked on maps, and the lines themselves marked with conspicuous plastic reflectors, to warn pilots of the presence of conductors.
Construction of overhead power lines, especially in wilderness areas, may have significant environmental effects. Environmental studies for such projects may consider the effect of bush clearing, changed migration routes for migratory animals, possible access by predators and humans along transmission corridors, disturbances of fish habitat at stream crossings, and other effects.
Aviation accidents
General aviation, hang gliding, paragliding, skydiving, balloon, and kite flying must avoid accidental contact with power lines. Nearly every kite product warns users to stay away from power lines. Deaths occur when aircraft crash into power lines. Some power lines are marked with obstruction markers, especially near air strips or over waterways that may support floatplane operations. The placement of power lines sometimes use up sites that would otherwise be used by hang gliders.
History
The first transmission of electrical impulses over an extended distance was demonstrated on July 14, 1729 by the physicist Stephen Gray. The demonstration used damp hemp cords suspended by silk threads (the low resistance of metallic conductors not being appreciated at the time).
However the first practical use of overhead lines was in the context of telegraphy. By 1837 experimental commercial telegraph systems ran as far as 20 km (13 miles). Electric power transmission was accomplished in 1882 with the first high-voltage transmission between Munich and Miesbach (60 km). 1891 saw the construction of the first three-phase alternating current overhead line on the occasion of the International Electricity Exhibition in Frankfurt, between Lauffen and Frankfurt.
In 1912 the first 110 kV-overhead power line entered service followed by the first 220 kV-overhead power line in 1923. In the 1920s RWE AG built the first overhead line for this voltage and in 1926 built a Rhine crossing with the pylons of Voerde, two masts 138 meters high.
In 1953, the first 345 kV line was put into service by American Electric Power in the United States. In Germany in 1957 the first 380 kV overhead power line was commissioned (between the transformer station and Rommerskirchen). In the same year the overhead line traversing of the Strait of Messina went into service in Italy, whose pylons served the Elbe crossing 1. This was used as the model for the building of the Elbe crossing 2 in the second half of the 1970s which saw the construction of the highest overhead line pylons of the world. Earlier, in 1952, the first 380 kV line was put into service in Sweden, in 1000 km (625 miles) between the more populated areas in the south and the largest hydroelectric power stations in the north.
Starting from 1967 in Russia, and also in the USA and Canada, overhead lines for voltage of 765 kV were built. In 1985 overhead power line was built in Soviet Union between Kokshetau and the power station at Ekibastuz, this was a three-phase alternating current line at 1150 kV. In 1999, in Japan the first powerline designed for 1000 kV with 2 circuits were built, the Kita-Iwaki Powerline. In 2002 the building of the highest overhead line commenced in China, the Yangtze River Crossing, its two high suspension towers beginning service in 2004.
In the 21st century, replacing steel with carbon fiber cores (advanced reconductoring) became a way for utilities to increase transmission capacity without increasing the amount of land used.
Mathematical analysis
An overhead power line is one example of a transmission line. At power system frequencies, many useful simplifications can be made for lines of typical lengths. For analysis of power systems, the distributed resistance, series inductance, shunt leakage resistance and shunt capacitance can be replaced with suitable lumped values or simplified networks.
Short and medium line model
A short length of a power line (less than 80 km) can be approximated with a resistance in series with an inductance and ignoring the shunt admittances. This value is not the total impedance of the line, but rather the series impedance per unit length of line. For a longer length of line (80–250 km), a shunt capacitance is added to the model. In this case it is common to distribute half of the total capacitance to each side of the line. As a result, the power line can be represented as a two-port network, such as with ABCD parameters.
The circuit can be characterized as
where
Z is the total series line impedance
z is the series impedance per unit length
l is the line length
is the sinusoidal angular frequency
The medium line has an additional shunt admittance
where
Y is the total shunt line admittance
y is the shunt admittance per unit length
| Technology | Electricity transmission and distribution | null |
1355078 | https://en.wikipedia.org/wiki/Overhead%20cable | Overhead cable | An overhead cable is a cable for the transmission of information, laid on utility poles. Overhead telephone and cable TV lines are common in North America. These poles sometimes carry overhead power lines for the supply of electric power. Power supply companies may also use them for an in-house communication network. Sometimes these cables are integrated in the ground or power conductor. Otherwise an additional line is strung on the masts.
Cables are arranged on poles with the most dangerous cables, that is, those carrying power, strung highest. Overhead cable systems also include a number of different components for managing signal cables. These include splicing systems that allow multi-conductor cables for distributing telephone signals and snowshoe-shaped devices for reversing the direction of cables.
When metal-based telephone wires are strung on the same utility poles as the power lines, they can pick up noise from the power line. Modern fiber optic telephone cable has the advantage that it can be strung next to power lines without interference.
In heavily populated regions of the UK, the only overhead cable that would be visible is the telephone line. Power cables and fiber-optic cables that deliver television and broadband services are buried underground. The lesser populated regions of the UK, the countryside for example, will have overhead power cables. Although it is safer to keep the cables underground, it would be difficult to repair a line if a fault were to develop.
| Technology | Electricity transmission and distribution | null |
1355101 | https://en.wikipedia.org/wiki/Feeler%20gauge | Feeler gauge | A feeler gauge is a tool used to measure gap widths. Feeler gauges are mostly used in engineering to measure the clearance between two parts.
Description
They consist of a number of small lengths of steel of different thicknesses with measurements marked on each piece. They are flexible enough that, even if they are all on the same hinge, several can be stacked together to gauge intermediate values. It is common to have two sets: one for imperial units (typically measured in thousandths of an inch), and one for metric (typically measured in hundredths of a millimetre) measurements (with intervals of thousandths of an inch and hundredths of a millimetre being roughly in the same order of magnitude).
The same device with wires of specific diameter instead of flat blades is used to set the gap in spark plugs to the correct size; this is done by increasing or decreasing the gap until the gauge of the correct size just fits inside the gap.
The lengths of steel are sometimes called leaves or blades, although they have no sharp edge.
Stainless steel is a common material for feeler gauges. Some feeler gauge sets have a single blade of brass due to the historical reason that early electronic ignition systems required the air gap between the reluctor and the pickup part being set with a non-ferrous metal.
Types
Taper
A taper feeler gauge is a feeler gauge of tapered, as opposed to parallel, shape. The blade of the gauge is of a constant thickness, and the two types of gauge are used in a similar way.
| Technology | Measuring instruments | null |
1356031 | https://en.wikipedia.org/wiki/Agribusiness | Agribusiness | Agribusiness is the industry, enterprises, and the field of study of value chains in agriculture and in the bio-economy,
in which case it is also called bio-business or bio-enterprise.
The primary goal of agribusiness is to maximize profit while satisfying the needs of consumers for products related to natural resources. Agribusinesses comprise farms, food and fiber processing, forestry, fisheries, biotechnology and biofuel enterprises and their input suppliers.
Studies of business growth and performance in farming have found that successful agricultural businesses are cost-efficient internally and operate in favourable economic, political, and physical-organic environments. They are able to expand and make profits, improve the productivity of land, labor, and capital, and keep their costs down to ensure market price competitiveness.
Agribusiness is not limited to farming. It encompasses a broader spectrum through the agribusiness system which includes input supplies, value-addition, marketing, entrepreneurship, microfinancing, and agricultural extension.
In some countries like the Philippines, creation and management of agribusiness enterprises require consultation with registered agriculturists above a certain level of operations, capitalization, land area, or number of animals in the farm.
Evolution of the agribusiness concept
The word "agribusiness" is a portmanteau of the words agriculture and business. The earliest known use of the word was in the Volume 155 of the Canadian Almanac & Directory published in 1847. Although most practitioners recognize that it was coined in 1957 by two Harvard Business School professors, John Davis and Ray Goldberg after they published the book "A Concept of Agribusiness."
Their book argued against the New Deal programs of then U.S. President Franklin Roosevelt as it led to the increase in agricultural prices. Davis and Goldberg favored corporate-driven agriculture or large-scale farming to revolutionize the agriculture sector, lessening the dependency on state power and politics. They explained in the book that vertically integrated firms within the agricultural value chains have the ability to control prices and where they are distributed. Goldberg then assisted in the establishment of the first undergraduate program in agribusiness in 1966 at the UP College of Agriculture in Los Baños, Philippines as Bachelor of Science in Agriculture major in Agribusiness. The program was initially a joint undertaking with the UP College of Business Administration in Diliman, Quezon City until 1975. Jose D. Drilon of the University of the Philippines then published the book "Agribusiness Management Resource Materials" (1971) which would be the foundation of current agribusiness programs around the world. In 1973, Drilon and Goldberg further expanded the concept of agribusiness to include support organizations such as governments, research institutions, schools, financial institutions, and cooperatives within the integrated Agribusiness System.
Mark R. Edwards and Clifford J. Shultz II (2005) of Loyola University Chicago reframed the definition of agribusiness to emphasize its lack of focus on farm production but towards market centricity and innovative approach to serve consumers worldwide.
In 2012, Thomas L. Sporleder and Michael A. Boland defined the unique economic characteristics of agribusiness supply chains from industrial manufacturing and service supply chains. They have identified seven main characteristics:
Risks emanating from the biological nature of agrifood supply chains
The role of buffer stocks within the supply chain
The scientific foundation of innovation in production agriculture having shifted from chemistry to biology
Cyberspace and information technology influences on agrifood supply chains
The prevalent market structure at the farm gate remains oligopsony
Relative market power shifts in agrifood supply chains away from food manufacturers downstream to food retailers
Globalization of agriculture and agrifood supply chains
In 2017, noting the rise of genetic engineering and biotechnology in agriculture, Goldberg further expanded the definition of agribusiness which covers all the interdependent aspects of the food system including medicine, nutrition, and health. He also emphasized the responsibility of agribusiness to be environmentally and socially conscious towards sustainability.
Some agribusinesses have adopted the triple bottom line framework such as aligning for fair trade, organic, good agricultural practices, and B-corporation certifications towards the concept of social entrepreneurship.
Agribusiness system
Inputs sector
Agricultural supplies
An agricultural supply store or agrocenter is an agriculturally-oriented shop where one sells agricultural supplies — inputs required for agricultural production such as pesticides, feed and fertilizers . Sometimes these stores are organized as cooperatives, where store customers aggregate their resources to purchase agricultural inputs. Agricultural supply and the stores that provide it are part of the larger Agribusiness industry.
Agricultural labor
Irrigation
Seeds
Fertilizers
Production sector
Farming
Farm mechanization
Processing sector
Primary processing
Secondary processing
Marketing sector
Farmers' market
Support sector
Education
Cooperatives
Governments
Professionals
Studies and reports
Studies of agribusiness often come from the academic fields of agricultural economics and management studies, sometimes called agribusiness management. To promote more development of food economies, many government agencies support the research and publication of economic studies and reports exploring agribusiness and agribusiness practices. Some of these studies are on foods produced for export and are derived from agencies focused on food exports. These agencies include the Foreign Agricultural Service (FAS) of the U.S. Department of Agriculture, Agriculture and Agri-Food Canada (AAFC), Austrade, and New Zealand Trade and Enterprise (NZTE).
The Federation of International Trade Associations publishes studies and reports by FAS and AAFC, as well as other non-governmental organizations on its website.
In their book A Concept of Agribusiness, Ray Goldberg and John Davis provided a rigorous economic framework for the field. They traced a complex value-added chain that begins with the farmer's purchase of seed and livestock and ends with a product fit for the consumer's table. Agribusiness boundary expansion is driven by a variety of transaction costs.
As concern over global warming intensifies, biofuels derived from crops are gaining increased public and scientific attention. This is driven by factors such as oil price spikes, the need for increased energy security, concern over greenhouse gas emissions from fossil fuels, and support from government subsidies. In Europe and in the US, increased research and production of biofuels have been mandated by law.
| Technology | Agriculture, labor and economy | null |
1356402 | https://en.wikipedia.org/wiki/Steric%20effects | Steric effects | Steric effects arise from the spatial arrangement of atoms. When atoms come close together there is generally a rise in the energy of the molecule. Steric effects are nonbonding interactions that influence the shape (conformation) and reactivity of ions and molecules. Steric effects complement electronic effects, which dictate the shape and reactivity of molecules. Steric repulsive forces between overlapping electron clouds result in structured groupings of molecules stabilized by the way that opposites attract and like charges repel.
Steric hindrance
Steric hindrance is a consequence of steric effects. Steric hindrance is the slowing of chemical reactions due to steric bulk. It is usually manifested in intermolecular reactions, whereas discussion of steric effects often focus on intramolecular interactions. Steric hindrance is often exploited to control selectivity, such as slowing unwanted side-reactions.
Steric hindrance between adjacent groups can also affect torsional bond angles. Steric hindrance is responsible for the observed shape of rotaxanes and the low rates of racemization of 2,2'-disubstituted biphenyl and binaphthyl derivatives.
Measures of steric properties
Because steric effects have profound impact on properties, the steric properties of substituents have been assessed by numerous methods.
Rate data
Relative rates of chemical reactions provide useful insights into the effects of the steric bulk of substituents. Under standard conditions, methyl bromide solvolyzes 107 faster than does neopentyl bromide. The difference reflects the inhibition of attack on the compound with the sterically bulky (CH3)3C group.
A-values
A-values provide another measure of the bulk of substituents. A-values are derived from equilibrium measurements of monosubstituted cyclohexanes. The extent that a substituent favors the equatorial position gives a measure of its bulk.
Ceiling temperatures
Ceiling temperature () is a measure of the steric properties of the monomers that comprise a polymer. is the temperature where the rate of polymerization and depolymerization are equal. Sterically hindered monomers give polymers with low 's, which are usually not useful.
Cone angles
Ligand cone angles are measures of the size of ligands in coordination chemistry. It is defined as the solid angle formed with the metal at the vertex and the hydrogen atoms at the perimeter of the cone (see figure).
Significance and applications
Steric effects are critical to chemistry, biochemistry, and pharmacology. In organic chemistry, steric effects are nearly universal and affect the rates and activation energies of most chemical reactions to varying degrees.
In biochemistry, steric effects are often exploited in naturally occurring molecules such as enzymes, where the catalytic site may be buried within a large protein structure. In pharmacology, steric effects determine how and at what rate a drug will interact with its target bio-molecules.
| Physical sciences | Stereochemistry | Chemistry |
1356480 | https://en.wikipedia.org/wiki/Lyon%20Metro | Lyon Metro | The Lyon Metro (, ) is a rapid transit system serving Lyon Metropolis, France. First opened in 1974, it currently consists of four lines, serving 42 stations and comprising of route. Part of the Transports en Commun Lyonnais (TCL) system of public transport, it is supported by two funiculars and a tramway network.
Unlike other French metro systems, but like RER and other SNCF services, Lyon Metro trains run on the left. This is the result of an unrealised project to run the metro into the suburbs on existing railway lines. The loading gauge for all lines is , more generous than the average for metros in Europe. The Lyon Metro has rubber-wheel cars. In 2018, the average daily weekday ridership was 755,000.
Routes
The Lyon Metro consists of four lines, A, B, C and D, each identified on maps by different colours:
Lines A and B
Line A from Perrache to Laurent Bonnevay–Astroballe and Line B from Charpennes to Part-Dieu were constructed by cut-and-cover and went into service on 2 May 1978, as the inaugural lines of the Lyon Metro. Trains on both lines run on rubber tyres rather than steel wheels.
Line B was extended to Jean Macé on 9 September 1981, to Stade de Gerland on 4 September 2000 as well as later to Gare d'Oullins on 11 December 2013.
An extension to Vaulx-en-Velin–La Soie on Line A opened in October 2007.
Since 2022, Line B is automated with new MPL 16 rolling stock ordered to Alstom in 2016. The MPL 75 trains previously used on Line B are meant to join the other MPL 75s on Line A to increase capacity.
An extension to Line B saw two stations, Oullins Centre and Saint-Genis-Laval–Hôpital Lyon Sud open on 20 October 2023.
Line C
The Croix-Rousse-Croix-Paquet rack railway, which was refurbished in 1974, was integrated into the Metro in 1978 as Line C, with an extension to Hôtel de Ville–Louis Pradel (thus running from Hôtel de Ville–Louis Pradel to Croix-Rousse). It was extended to Cuire on 8 December 1984.
The line was constructed using various methods; the incline rising through a deep tunnel, the portion on the flat at Croix-Rousse using cut-and-cover while the section beyond Hénon runs on the surface. The Croix Paquet station claims to be the steepest metro station in Europe, with an incline of 17%.
Line C uses overhead wires and steel wheels while Lines A, B and D use a third rail and rubber tyres. Until Paris Métro Line 15 opens it is the only metro line in France to use overhead lines and the only steel wheeled metro line in France outside Paris.
Line D
Line D, the first fully automatic metro line in France, started with operators on board trains on 4 September 1991, between Gorge de Loup and Grange Blanche. The line was extended to Gare de Vénissieux on 11 December 1992, when it switched to driverless operation. On 28 April 1997, it was extended again to Gare de Vaise.
Using rubber tyres like lines A and B, trains on line D are controlled by a system known as MAGGALY (Métro Automatique à Grand Gabarit de l’Agglomération Lyonnaise). Unusually for a driverless metro, no platform screen doors are installed on station platforms. The trains use infrared sensors to detect obstructions on the track. Other systems using this technology include the Nuremberg U-Bahn and Budapest Metro's Line 4.
The deepest line in Lyon, Line D was constructed partly using boring machines and passes under both rivers, the Rhône and Saône. At long with 15 stations, it is also the longest line in Lyon.
In 2016, new MPL 16 rolling stock was ordered from Alstom for Line B and Line D; it came into service on Line B in 2022. These trains allow for an increase in capacity on Line D. Further, they will be coupled to form four-car units at rush hours and should replace the MPL 75 of Line B which would then solely run on Line A.
Map
Operation
The Metro, like the rest of the local public transport system, is operated by Keolis Lyon (ex-SLTC - the (Lyon public transport company)), under the TCL brand - (Lyon public transport). It is operated on behalf of SYTRAL Mobilités - the (Rhône department and Lyon metropolitan transport syndicate), a Syndicat Mixte. On 1 January 2025 RATP Dev will take over operation of the metro.
Future expansion
A new line, dubbed Line E, was under consideration to link Lyon's western suburbs to the city centre. Twelve variants were initially proposed; two options, running from either Bellecour or Hôtel de Ville to Alaï, were selected for further study and could potentially have been opened around 2030. In 2022, however, the plans for Line E and other metro extensions have been cancelled in favor of plans for new express tramways, partly underground.
| Technology | France | null |
1357154 | https://en.wikipedia.org/wiki/Parallel%20evolution | Parallel evolution | Parallel evolution is the similar development of a trait in distinct species that are not closely related, but share a similar original trait in response to similar evolutionary pressure.
Parallel vs. convergent evolution
Given a trait that occurs in each of two lineages descended from a specified ancestor, it is possible in theory to define parallel and convergent evolutionary trends strictly, and distinguish them clearly from one another. However, the criteria for defining convergent as opposed to parallel evolution are unclear in practice, so that arbitrary diagnosis is common. When two species share a trait, evolution is defined as parallel if the ancestors are known to have shared that similarity; if not, it is defined as convergent. However, the stated conditions are a matter of degree; all organisms share common ancestors. Scientists differ on whether the distinction is useful.
Parallel evolution between marsupials and placentals
A number of examples of parallel evolution are provided by the two main branches of the mammals, the placentals and marsupials, which have followed independent evolutionary pathways following the break-up of land-masses such as Gondwanaland roughly 100 million years ago. In South America, marsupials and placentals shared the ecosystem (before the Great American Interchange); in Australia, marsupials prevailed; and in the Old World and North America the placentals won out. However, in all these localities mammals were small and filled only limited places in the ecosystem until the mass extinction of dinosaurs sixty-five million years ago. At this time, mammals on all three landmasses began to take on a much wider variety of forms and roles. While some forms were unique to each environment, surprisingly similar animals have often emerged in two or three of the separated continents. Examples of these include the placental sabre-toothed cats (Machairodontinae) and the South American marsupial sabre-tooth (Thylacosmilus); the Tasmanian wolf and the European wolf; likewise marsupial and placental moles, flying squirrels, and (arguably) mice.
Parallel coevolution of traits between hummingbirds and sunbirds contributing to ecological guilds
Hummingbirds and sunbirds, two nectarivorous bird lineages in the New and Old Worlds have parallelly evolved a suite of specialized behavioral and anatomical traits. These traits (bill shape, digestive enzymes, and flight) allow the birds to optimally fit the flower-feeding-and-pollination ecological niche they occupy, which is shaped by the birds' suites of parallel traits. Thus, a parallel coevolved behavioral syndrome within the birds creates an emergent guild of highly specialized birds and highly adapted plants, each exploiting the other's involvement in the flowers' pollination in the Old World and New World alike.
The bill shape of nectarivores, being long and needle-like, allows them to reach down a flower's pistil/stamen and get at the nectar within. Nectarivores may also use their specialized bills to engage in nectar robbing, a practice seen in both hummingbirds and sunbirds in which the bird gets nectar by making a hole in the base of the flower's corolla tube instead of inserting its bill through the tube as is standard, thus "robbing" the flower of nectar since it is not pollinated it in return.
Nectarivores and ornithophilous flowers often exist in mutualistic guild relationships facilitated by the bird's bill shape, food source, and digestive ability acting in concert with the flower's tube shape and adaptation to pollination by hovering or perching birds. The birds eat nectar using their long, thin bills and, in so doing, collect pollen on their bills; this pollen is then transferred to the next flower they feed on. This mutualism coevolved in parallel between the Old World and New World birds and their respective flowers. Moreover, the digestive enzyme activity in nectarivores matching the nectar composition in their respective flowers appears to have coevolved in parallel between plants and pollinators across continents, as the nectarivorous lineages independently evolved the ability to digest the nectar specific to their flowers, resulting in distinct guilds.
The capacity of nectarivores to digest sucrose is far greater than that of other avian taxa. This difference is due to an analogous high concentration of sucrase-isomaltase, an enzyme that hydrolyzes sucrose. Sucrase activity per unit intestinal surface area appears to be higher in nectarivores than in other birds, meaning these nectarivorous avians can digest more sucrose more rapidly than other taxa. Moreover, the Adaptive Modulation Hypothesis does not apply for nectarivores and sugar-digesting enzymes, meaning that two lineages of nectarivores should not necessarily both have high sucrase-isomaltase concentrations even though they both eat nectar. Thus, parallel acquisition of analogous sucrose digestive capability is a reasonable conclusion because there is no apparent cause for the two lineages to share this high enzyme concentration.
| Biology and health sciences | Basics_4 | Biology |
1357514 | https://en.wikipedia.org/wiki/Recombinant%20DNA | Recombinant DNA | Recombinant DNA (rDNA) molecules are DNA molecules formed by laboratory methods of genetic recombination (such as molecular cloning) that bring together genetic material from multiple sources, creating sequences that would not otherwise be found in the genome.
Recombinant DNA is the general name for a piece of DNA that has been created by combining two or more fragments from different sources. Recombinant DNA is possible because DNA molecules from all organisms share the same chemical structure, differing only in the nucleotide sequence. Recombinant DNA molecules are sometimes called chimeric DNA because they can be made of material from two different species like the mythical chimera. rDNA technology uses palindromic sequences and leads to the production of sticky and blunt ends.
The DNA sequences used in the construction of recombinant DNA molecules can originate from any species. For example, plant DNA can be joined to bacterial DNA, or human DNA can be joined with fungal DNA. In addition, DNA sequences that do not occur anywhere in nature can be created by the chemical synthesis of DNA and incorporated into recombinant DNA molecules. Using recombinant DNA technology and synthetic DNA, any DNA sequence can be created and introduced into living organisms.
Proteins that can result from the expression of recombinant DNA within living cells are termed recombinant proteins. When recombinant DNA encoding a protein is introduced into a host organism, the recombinant protein is not necessarily produced. Expression of foreign proteins requires the use of specialized expression vectors and often necessitates significant restructuring by
foreign coding sequences.
Recombinant DNA differs from genetic recombination in that the former results from artificial methods while the latter is a normal biological process that results in the remixing of existing DNA sequences in essentially all organisms.
Production
Molecular cloning is the laboratory process used to produce recombinant DNA. It is one of two most widely used methods, along with polymerase chain reaction (PCR), used to direct the replication of any specific DNA sequence chosen by the experimentalist. There are two fundamental differences between the methods. One is that molecular cloning involves replication of the DNA within a living cell, while PCR replicates DNA in the test tube, free of living cells. The other difference is that cloning involves cutting and pasting DNA sequences, while PCR amplifies by copying an existing sequence.
Formation of recombinant DNA requires a cloning vector, a DNA molecule that replicates within a living cell. Vectors are generally derived from plasmids or viruses, and represent relatively small segments of DNA that contain necessary genetic signals for replication, as well as additional elements for convenience in inserting foreign DNA, identifying cells that contain recombinant DNA, and, where appropriate, expressing the foreign DNA. The choice of vector for molecular cloning depends on the choice of host organism, the size of the DNA to be cloned, and whether and how the foreign DNA is to be expressed. The DNA segments can be combined by using a variety of methods, such as restriction enzyme/ligase cloning or Gibson assembly.
In standard cloning protocols, the cloning of any DNA fragment essentially involves seven steps: (1) Choice of host organism and cloning vector, (2) Preparation of vector DNA, (3) Preparation of DNA to be cloned, (4) Creation of recombinant DNA, (5) Introduction of recombinant DNA into the host organism, (6) Selection of organisms containing recombinant DNA, and (7) Screening for clones with desired DNA inserts and biological properties.
These steps are described in some detail in a related article (molecular cloning).
DNA expression
DNA expression requires the transfection of suitable host cells. Typically, either bacterial, yeast, insect, or mammalian cells (such as Human Embryonic Kidney cells or CHO cells) are used as host cells.
Following transplantation into the host organism, the foreign DNA contained within the recombinant DNA construct may or may not be expressed. That is, the DNA may simply be replicated without expression, or it may be transcribed and translated and a recombinant protein is produced. Generally speaking, expression of a foreign gene requires restructuring the gene to include sequences that are required for producing an mRNA molecule that can be used by the host's translational apparatus (e.g. promoter, translational initiation signal, and transcriptional terminator). Specific changes to the host organism may be made to improve expression of the ectopic gene. In addition, changes may be needed to the coding sequences as well, to optimize translation, make the protein soluble, direct the recombinant protein to the proper cellular or extracellular location, and stabilize the protein from degradation.
Properties of organisms containing recombinant DNA
In most cases, organisms containing recombinant DNA have apparently normal phenotypes. That is, their appearance, behavior and metabolism are usually unchanged, and the only way to demonstrate the presence of recombinant sequences is to examine the DNA itself, typically using a polymerase chain reaction (PCR) test. Significant exceptions exist, and are discussed below.
If the rDNA sequences encode a gene that is expressed, then the presence of RNA and/or protein products of the recombinant gene can be detected, typically using RT-PCR or western hybridization methods. Gross phenotypic changes are not the norm, unless the recombinant gene has been chosen and modified so as to generate biological activity in the host organism. Additional phenotypes that are encountered include toxicity to the host organism induced by the recombinant gene product, especially if it is over-expressed or expressed within inappropriate cells or tissues.
In some cases, recombinant DNA can have deleterious effects even if it is not expressed. One mechanism by which this happens is insertional inactivation, in which the rDNA becomes inserted into a host cell's gene. In some cases, researchers use this phenomenon to "knock out" genes to determine their biological function and importance. Another mechanism by which rDNA insertion into chromosomal DNA can affect gene expression is by inappropriate activation of previously unexpressed host cell genes. This can happen, for example, when a recombinant DNA fragment containing an active promoter becomes located next to a previously silent host cell gene, or when a host cell gene that functions to restrain gene expression undergoes insertional inactivation by recombinant DNA.
Applications of recombinant DNA
Recombinant DNA is widely used in biotechnology, medicine and research. Today, recombinant proteins and other products that result from the use of DNA technology are found in essentially every pharmacy, physician or veterinarian office, medical testing laboratory, and biological research laboratory. In addition, organisms that have been manipulated using recombinant DNA technology, as well as products derived from those organisms, have found their way into many farms, supermarkets, home medicine cabinets, and even pet shops, such as those that sell GloFish and other genetically modified animals.
The most common application of recombinant DNA is in basic research, in which the technology is important to most current work in the biological and biomedical sciences. Recombinant DNA is used to identify, map and sequence genes, and to determine their function. rDNA probes are employed in analyzing gene expression within individual cells, and throughout the tissues of whole organisms. Recombinant proteins are widely used as reagents in laboratory experiments and to generate antibody probes for examining protein synthesis within cells and organisms.
Many additional practical applications of recombinant DNA are found in industry, food production, human and veterinary medicine, agriculture, and bioengineering. Some specific examples are identified below.
Recombinant chymosin
Found in rennet, chymosin is the enzyme responsible for hydrolysis of κ-casein to produce para-κ-casein and glycomacropeptide, which is the first step in formation of cheese, and subsequently curd, and whey. It was the first genetically engineered food additive used commercially. Traditionally, processors obtained chymosin from rennet, a preparation derived from the fourth stomach of milk-fed calves. Scientists engineered a non-pathogenic strain (K-12) of E. coli bacteria for large-scale laboratory production of the enzyme. This microbiologically produced recombinant enzyme, identical structurally to the calf derived enzyme, costs less and is produced in abundant quantities. Today about 60% of U.S. hard cheese is made with genetically engineered chymosin. In 1990, FDA granted chymosin "generally recognized as safe" (GRAS) status based on data showing that the enzyme was safe.
Recombinant human insulin
Recombinant human insulin has almost completely replaced insulin obtained from animal sources (e.g. pigs and cattle) for the treatment of type 1 diabetes. A variety of different recombinant insulin preparations are in widespread use. Recombinant insulin is synthesized by inserting the human insulin gene into E. coli, or yeast (Saccharomyces cerevisiae) which then produces insulin for human use. Insulin produced by E. coli requires further post translational modifications (e.g. glycosylation) whereas yeasts are able to perform these modifications themselves by virtue of being more complex host organisms. The advantage of recombinant human insulin is after chronic use patients don't develop an immune defence against it the way animal sourced insulin stimulates the human immune system.
Recombinant human growth hormone (HGH, somatotropin)
Administered to patients whose pituitary glands generate insufficient quantities to support normal growth and development. Before recombinant HGH became available, HGH for therapeutic use was obtained from pituitary glands of cadavers. This unsafe practice led to some patients developing Creutzfeldt–Jakob disease. Recombinant HGH eliminated this problem, and is now used therapeutically. It has also been misused as a performance-enhancing drug by athletes and others.
Recombinant blood clotting factor VIII
It is the recombinant form of factor VIII, a blood-clotting protein that is administered to patients with the bleeding disorder hemophilia, who are unable to produce factor VIII in quantities sufficient to support normal blood coagulation. Before the development of recombinant factor VIII, the protein was obtained by processing large quantities of human blood from multiple donors, which carried a very high risk of transmission of blood borne infectious diseases, for example HIV and hepatitis B.
Recombinant hepatitis B vaccine
Hepatitis B infection can be successfully controlled through the use of a recombinant subunit hepatitis B vaccine, which contains a form of the hepatitis B virus surface antigen that is produced in yeast cells. The development of the recombinant subunit vaccine was an important and necessary development because hepatitis B virus, unlike other common viruses such as polio virus, cannot be grown in vitro.
Recombinant antibodies
Recombinant antibodies (rAbs) are produced in vitro by the means of expression systems based on mammalian cells. Their monospecific binding to a specific epitope makes rAbs eligible not only for research purposes, but also as therapy options against certain cancer types, infections and autoimmune diseases.
Diagnosis of HIV infection
Each of the three widely used methods for diagnosing HIV infection has been developed using recombinant DNA. The antibody test (ELISA or western blot) uses a recombinant HIV protein to test for the presence of antibodies that the body has produced in response to an HIV infection. The DNA test looks for the presence of HIV genetic material using reverse transcription polymerase chain reaction (RT-PCR). Development of the RT-PCR test was made possible by the molecular cloning and sequence analysis of HIV genomes. HIV testing page from US Centers for Disease Control (CDC)
Golden rice
Golden rice is a recombinant variety of rice that has been engineered to express the enzymes responsible for β-carotene biosynthesis. This variety of rice holds substantial promise for reducing the incidence of vitamin A deficiency in the world's population. Golden rice is not currently in use, pending the resolution of regulatory and intellectual property issues.
Herbicide-resistant crops
Commercial varieties of important agricultural crops (including soy, maize/corn, sorghum, canola, alfalfa and cotton) have been developed that incorporate a recombinant gene that results in resistance to the herbicide glyphosate (trade name Roundup), and simplifies weed control by glyphosate application. These crops are in common commercial use in several countries.
Insect-resistant crops
Bacillus thuringiensis is a bacterium that naturally produces a protein (Bt toxin) with insecticidal properties. The bacterium has been applied to crops as an insect-control strategy for many years, and this practice has been widely adopted in agriculture and gardening. Recently, plants have been developed that express a recombinant form of the bacterial protein, which may effectively control some insect predators. Environmental issues associated with the use of these transgenic crops have not been fully resolved.
History
The idea of recombinant DNA was first proposed by Peter Lobban, a graduate student of Prof. Dale Kaiser in the Biochemistry Department at Stanford University Medical School. The first publications describing the successful production and intracellular replication of recombinant DNA appeared in 1972 and 1973, from Stanford and UCSF. In 1980 Paul Berg, a professor in the Biochemistry Department at Stanford and an author on one of the first papers was awarded the Nobel Prize in Chemistry for his work on nucleic acids "with particular regard to recombinant DNA". Werner Arber, Hamilton Smith, and Daniel Nathans shared the 1978 Nobel Prize in Physiology or Medicine for the discovery of restriction endonucleases which enhanced the techniques of rDNA technology.
Stanford University applied for a U.S. patent on recombinant DNA on November 4, 1974, listing the inventors as Herbert W. Boyer (professor at the University of California, San Francisco) and Stanley N. Cohen (professor at Stanford University); this patent, U.S. 4,237,224A, was awarded on December 2, 1980. The first licensed drug generated using recombinant DNA technology was human insulin, developed by Genentech and licensed by Eli Lilly and Company.
Controversy
Scientists associated with the initial development of recombinant DNA methods recognized that the potential existed for organisms containing recombinant DNA to have undesirable or dangerous properties. At the 1975 Asilomar Conference on Recombinant DNA, these concerns were discussed and a voluntary moratorium on recombinant DNA research was initiated for experiments that were considered particularly risky. This moratorium was widely observed until the US National Institutes of Health developed and issued formal guidelines for rDNA work. Today, recombinant DNA molecules and recombinant proteins are usually not regarded as dangerous. However, concerns remain about some organisms that express recombinant DNA, particularly when they leave the laboratory and are introduced into the environment or food chain. These concerns are discussed in the articles on genetically modified organisms and genetically modified food controversies. Furthermore, there are concerns about the by-products in biopharmaceutical production, where recombinant DNA result in specific protein products. The major by-product, termed host cell protein, comes from the host expression system and poses a threat to the patient's health and the overall environment.
| Technology | Biotechnology | null |
1358431 | https://en.wikipedia.org/wiki/Surface%20brightness | Surface brightness | In astronomy, surface brightness (SB) quantifies the apparent brightness or flux density per unit angular area of a spatially extended object such as a galaxy or nebula, or of the night sky background. An object's surface brightness depends on its surface luminosity density, i.e., its luminosity emitted per unit surface area. In visible and infrared astronomy, surface brightness is often quoted on a magnitude scale, in magnitudes per square arcsecond (MPSAS) in a particular filter band or photometric system.
Measurement of the surface brightnesses of celestial objects is called surface photometry.
General description
The total magnitude is a measure of the brightness of an extended object such as a nebula, cluster, galaxy or comet. It can be obtained by summing up the luminosity over the area of the object. Alternatively, a photometer can be used by applying apertures or slits of different sizes of diameter. The background light is then subtracted from the measurement to obtain the total brightness. The resulting magnitude value is the same as a point-like source that is emitting the same amount of energy. The total magnitude of a comet is the combined magnitude of the coma and nucleus.
The apparent magnitude of an astronomical object is generally given as an integrated value—if a galaxy is quoted as having a magnitude of 12.5, it means we see the same total amount of light from the galaxy as we would from a star with magnitude 12.5. However, a star is so small it is effectively a point source in most observations (the largest angular diameter, that of R Doradus, is 0.057 ± 0.005 arcsec), whereas a galaxy may extend over several arcseconds or arcminutes. Therefore, the galaxy will be harder to see than the star against the airglow background light. Apparent magnitude is a good indication of visibility if the object is point-like or small, whereas surface brightness is a better indicator if the object is large. What counts as small or large depends on the specific viewing conditions and follows from Ricco's law. In general, in order to adequately assess an object's visibility one needs to know both parameters.
This is the reason the extreme naked eye limit for viewing a star is apparent magnitude 8, but only apparent magnitude 6.9 for galaxies.
Calculating surface brightness
Surface brightnesses are usually quoted in magnitudes per square arcsecond. Because the magnitude is logarithmic, calculating surface brightness cannot be done by simple division of magnitude by area. Instead, for a source with a total or integrated magnitude m extending over a visual area of A square arcseconds, the surface brightness S is given by
For astronomical objects, surface brightness is analogous to photometric luminance and is therefore constant with distance: as an object becomes fainter with distance, it also becomes correspondingly smaller in visual area. In geometrical terms, for a nearby object emitting a given amount of light, radiative flux decreases with the square of the distance to the object, but the physical area corresponding to a given solid angle or visual area (e.g. 1 square arcsecond) decreases by the same proportion, resulting in the same surface brightness. For extended objects such as nebulae or galaxies, this allows the estimation of spatial distance from surface brightness by means of the distance modulus or luminosity distance.
Relationship to physical units
The surface brightness in magnitude units is related to the surface brightness in physical units of solar luminosity per square parsec by
where and are the absolute magnitude and the luminosity of the Sun in chosen color-band respectively.
Surface brightness can also be expressed in candela per square metre using the formula [value in cd/m2] = × 10(−0.4×[value in mag/arcsec2]).
Examples
A truly dark sky has a surface brightness of cd m−2 or 21.8 mag arcsec−2.
The peak surface brightness of the central region of the Orion Nebula is about 17 Mag/arcsec2 (about 14 millinits) and the outer bluish glow has a peak surface brightness of 21.3 Mag/arcsec2 (about 0.27 millinits).
| Physical sciences | Basics | Astronomy |
1358453 | https://en.wikipedia.org/wiki/Astrophysical%20jet | Astrophysical jet | An astrophysical jet is an astronomical phenomenon where outflows of ionised matter are emitted as extended beams along the axis of rotation. When this greatly accelerated matter in the beam approaches the speed of light, astrophysical jets become relativistic jets as they show effects from special relativity.
The formation and powering of astrophysical jets are highly complex phenomena that are associated with many types of high-energy astronomical sources. They likely arise from dynamic interactions within accretion disks, whose active processes are commonly connected with compact central objects such as black holes, neutron stars or pulsars. One explanation is that tangled magnetic fields are organised to aim two diametrically opposing beams away from the central source by angles only several degrees wide Jets may also be influenced by a general relativity effect known as frame-dragging.
Most of the largest and most active jets are created by supermassive black holes (SMBH) in the centre of active galaxies such as quasars and radio galaxies or within galaxy clusters. Such jets can exceed millions of parsecs in length. Other astronomical objects that contain jets include cataclysmic variable stars, X-ray binaries and gamma-ray bursts (GRB). Jets on a much smaller scale (~parsecs) may be found in star forming regions, including T Tauri stars and Herbig–Haro objects; these objects are partially formed by the interaction of jets with the interstellar medium. Bipolar outflows may also be associated with protostars, or with evolved post-AGB stars, planetary nebulae and bipolar nebulae.
Relativistic jets
Relativistic jets are beams of ionised matter accelerated close to the speed of light. Most have been observationally associated with central black holes of some active galaxies, radio galaxies or quasars, and also by galactic stellar black holes, neutron stars or pulsars. Beam lengths may extend between several thousand, hundreds of thousands or millions of parsecs. Jet velocities when approaching the speed of light show significant effects of the special theory of relativity; for example, relativistic beaming that changes the apparent beam brightness.
Massive central black holes in galaxies have the most powerful jets, but their structure and behaviours are similar to those of smaller galactic neutron stars and black holes. These SMBH systems are often called microquasars and show a large range of velocities. SS 433 jet, for example, has a mean velocity of 0.26c. Relativistic jet formation may also explain observed gamma-ray bursts, which have the most relativistic jets known, being ultrarelativistic.
Mechanisms behind the composition of jets remain uncertain, though some studies favour models where jets are composed of an electrically neutral mixture of nuclei, electrons, and positrons, while others are consistent with jets composed of positron–electron plasma. Trace nuclei swept up in a relativistic positron–electron jet would be expected to have extremely high energy, as these heavier nuclei should attain velocity equal to the positron and electron velocity.
Rotation as possible energy source
Because of the enormous amount of energy needed to launch a relativistic jet, some jets are possibly powered by spinning black holes. However, the frequency of high-energy astrophysical sources with jets suggests combinations of different mechanisms indirectly identified with the energy within the associated accretion disk and X-ray emissions from the generating source. Two early theories have been used to explain how energy can be transferred from a black hole into an astrophysical jet:
Blandford–Znajek process. This theory explains the extraction of energy from magnetic fields around an accretion disk, which are dragged and twisted by the spin of the black hole. Relativistic material is then feasibly launched by the tightening of the field lines.
Penrose mechanism. Here energy is extracted from a rotating black hole by frame dragging, which was later theoretically proven by Reva Kay Williams to be able to extract relativistic particle energy and momentum, and subsequently shown to be a possible mechanism for jet formation. This effect includes using general relativistic gravitomagnetism.
Relativistic jets from neutron stars
Jets may also be observed from spinning neutron stars. An example is pulsar IGR J11014-6103, which has the largest jet so far observed in the Milky Way, and whose velocity is estimated at 80% the speed of light (0.8c). X-ray observations have been obtained, but there is no detected radio signature nor accretion disk. Initially, this pulsar was presumed to be rapidly spinning, but later measurements indicate the spin rate is only 15.9 Hz. Such a slow spin rate and lack of accretion material suggest the jet is neither rotation nor accretion powered, though it appears aligned with the pulsar rotation axis and perpendicular to the pulsar's true motion.
Other images
| Physical sciences | Basics_2 | Astronomy |
1358820 | https://en.wikipedia.org/wiki/Pavilion | Pavilion | In architecture, pavilion has several meanings;
It may be a subsidiary building that is either positioned separately or as an attachment to a main building. Often it is associated with pleasure. In palaces and traditional mansions of Asia, there may be pavilions that are either freestanding or connected by covered walkways, as in the Forbidden City (Chinese pavilions), Topkapi Palace in Istanbul, and in Mughal buildings like the Red Fort.
As part of a large palace, pavilions may be symmetrically placed building blocks that flank (appear to join) a main building block or the outer ends of wings extending from both sides of a central building block, the corps de logis. Such configurations provide an emphatic visual termination to the composition of a large building, akin to bookends.
The word is from French (Old French ) and it meant a small palace, from Latin (accusative of ). In Late Latin and Old French, it meant both ‘butterfly’ and ‘tent’, because the canvas of a tent resembled a butterfly's spread wings.
The word is from the early 13c., paviloun, "large, stately tent raised on posts and used as a movable habitation," from Old French paveillon "large tent; butterfly" (12c.), from Latin papilionem (nominative papilio) "butterfly, moth," in Medieval Latin "tent" (see papillon); the type of tent was so called on its resemblance to wings. Meaning "open building in a park, etc., used for shelter or entertainment" is attested from 1680s. Sense of "small or moderate-sized building, isolated from but dependent on a larger or principal building" (as in a hospital) is by 1858.
Free-standing structures
Pavilions may be small garden outbuildings, similar to a summer house or a kiosk; small rooms on the roof of a large house, reached only via the roof (rather than by internal stairs) may also be called pavilions. These were particularly popular up to the 18th century and can be equated to the Italian , formerly rendered in English "casino". These often resembled small classical temples and follies. Especially if there is some space for food preparation, they may be called a banqueting house. A pavilion built to take advantage of a view may be referred to as a gazebo. Bandstands in a park are a class of pavilion. A by a swimming pool may have sufficient character and charm to be called a pavilion. By contrast, a free-standing pavilion can also be a far larger building such as the Royal Pavilion at Brighton, which is in fact a large Indian-style palace; however, like its smaller namesakes, the common factor is that it was built for pleasure and relaxation.
A sports pavilion is usually a building adjacent to a sports ground used for changing clothes and often partaking of refreshments. Often it has a verandah to provide protection from the sun for spectators. In cricket grounds, as at Lord's, a cricket pavilion tends to be used for the building the players emerge from and return to, even when this is actually a large building including a grandstand. A pavilion in stadia, especially baseball parks, is a typically single-decked covered seating area (as opposed to the more expensive seating area of the main grandstand and the less expensive seating area of the uncovered bleachers).
Classical architecture
Externally, pavilions may be emphasised by any combination of a change in height, profile (a flat facade may end in round pavilions, or flat ones that project out), colour, material, and ornament. Internally they may be part of a rectangular block, or only connected to the main block by a thin section of building. The two 18th-century English country houses of Houghton Hall and Holkham Hall illustrate these different approaches in turn.
In the Place des Vosges (1605–1612), Paris, twin pavilions mark the centers of the north and south sides of the square. They are named the (“king’s pavilion”) and the (“queen’s pavilion”), though no royal personage ever lived in the square. With their triple archways, they function like gatehouses that give access to the privileged space of the square. French gatehouses had been built in the form of such pavilions in the preceding century.
Other uses
In some areas, a pavilion is a term for a hunting lodge. The in Luberon, France, is a typical 18th-century aristocratic hunting pavilion. The pavilion, located on the site of an old Roman villa, includes a garden , which was used by the guests for receptions.
Gallery
| Technology | Buildings and infrastructure | null |
18938226 | https://en.wikipedia.org/wiki/Digital%20rights%20management | Digital rights management | Digital rights management (DRM) is the management of legal access to digital content. Various tools or technological protection measures (TPM), such as access control technologies, can restrict the use of proprietary hardware and copyrighted works. DRM technologies govern the use, modification and distribution of copyrighted works (e.g. software, multimedia content) and of systems that enforce these policies within devices. DRM technologies include licensing agreements and encryption.
Laws in many countries criminalize the circumvention of DRM, communication about such circumvention, and the creation and distribution of tools used for such circumvention. Such laws are part of the United States' Digital Millennium Copyright Act (DMCA), and the European Union's Information Society Directive – with the French DADVSI an example of a member state of the European Union implementing that directive.
Copyright holders argue that DRM technologies are necessary to protect intellectual property, just as physical locks prevent personal property from theft. For examples, they can help the copyright holders for maintaining artistic controls, and supporting licenses' modalities such as rentals. Industrial users (i.e. industries) have expanded the use of DRM technologies to various hardware products, such as Keurig's coffeemakers, Philips' light bulbs, mobile device power chargers, and John Deere's tractors. For instance, tractor companies try to prevent farmers from making repairs via DRM.
DRM is controversial. There is an absence of evidence about the DRM capability in preventing copyright infringement, some complaints by legitimate customers for caused inconveniences, and a suspicion of stifling innovation and competition. Furthermore, works can become permanently inaccessible if the DRM scheme changes or if a required service is discontinued. DRM technologies have been criticized for restricting individuals from copying or using the content legally, such as by fair use or by making backup copies. DRM is in common use by the entertainment industry (e.g., audio and video publishers). Many online stores such as OverDrive use DRM technologies, as do cable and satellite service operators. Apple removed DRM technology from iTunes around 2009. Typical DRM also prevents lending materials out through a library, or accessing works in the public domain.
Introduction
The rise of digital media and analog-to-digital conversion technologies has increased the concerns of copyright-owners, particularly within the music and video industries. While analog media inevitably lose quality with each copy generation and during normal use, digital media files may be duplicated without limit with no degradation. Digital devices make it convenient for consumers to convert (rip) media originally in a physical, analog or broadcast form into a digital form for portability or later use. Combined with the Internet and file-sharing tools, made unauthorized distribution of copyrighted content (digital piracy) much easier.
History
DRM became a major concern with the growth of the Internet in the 1990s, as piracy crushed CD sales and online video became popular. It peaked in the early 2000s as various countries attempted to respond with legislation and regulations and dissipated in the 2010s as social media and streaming services largely replaced piracy and content providers elaborated next-generation business models.
Early efforts
In 1983, the Software Service System (SSS) devised by the Japanese engineer Ryuichi Moriya was the first example of DRM technology. It was subsequently refined under the name superdistribution. The SSS was based on encryption, with specialized hardware that controlled decryption and enabled payments to be sent to the copyright holder. The underlying principle was that the physical distribution of encrypted digital products should be completely unrestricted and that users of those products would be encouraged to do so.
An early DRM protection method for computer and Nintendo Entertainment System games was when the game would pause and prompt the player to look up a certain page in a booklet or manual that came with the game; if the player lacked access to the material, they would not be able to continue.
An early example of a DRM system is the Content Scramble System (CSS) employed by the DVD Forum on DVD movies. CSS uses an encryption algorithm to encrypt content on the DVD disc. Manufacturers of DVD players must license this technology and implement it in their devices so that they can decrypt the content. The CSS license agreement includes restrictions on how the DVD content is played, including what outputs are permitted and how such permitted outputs are made available. This keeps the encryption intact as the content is displayed.
In May 1998, the Digital Millennium Copyright Act (DMCA) passed as an amendment to US copyright law. It had controversial (possibly unintended) implications. Russian programmer Dmitry Sklyarov was arrested for alleged DMCA infringement after a presentation at DEF CON. The DMCA has been cited as chilling to legitimate users; such as security consultants including Niels Ferguson, who declined to publish vulnerabilities he discovered in Intel's secure-computing scheme due to fear of arrest under DMCA; and blind or visually impaired users of screen readers or other assistive technologies.
In 1999, Jon Lech Johansen released DeCSS, which allowed a CSS-encrypted DVD to play on a computer running Linux, at a time when no compliant DVD player for Linux had yet been created. The legality of DeCSS is questionable: one of its authors was sued, and reproduction of the keys themselves is subject to restrictions as illegal numbers.
More modern examples include ADEPT, FairPlay, Advanced Access Content System.
The World Intellectual Property Organization Copyright Treaty (WCT) was passed in 1996. The US Digital Millennium Copyright Act (DMCA), was passed in 1998. The European Union enacted the Information Society Directive. In 2006, the lower house of the French parliament adopted such legislation as part of the controversial DADVSI law, but added that protected DRM techniques should be made interoperable, a move which caused widespread controversy in the United States. The Tribunal de grande instance de Paris concluded in 2006, that the complete blocking of any possibilities of making private copies was an impermissible behaviour under French copyright law.
2000s
The broadcast flag concept was developed by Fox Broadcasting in 2001, and was supported by the MPAA and the U.S. Federal Communications Commission (FCC). A ruling in May 2005 by a United States courts of appeals held that the FCC lacked authority to impose it on the US TV industry. It required that all HDTVs obey a stream specification determining whether a stream can be recorded. This could block instances of fair use, such as time-shifting. It achieved more success elsewhere when it was adopted by the Digital Video Broadcasting Project (DVB), a consortium of about 250 broadcasters, manufacturers, network operators, software developers, and regulatory bodies from about 35 countries involved in attempting to develop new digital TV standards.
In January 2001, the Workshop on Digital Rights Management of the World Wide Web Consortium was held.
On 22 May 2001, the European Union passed the Information Society Directive, with copyright protections.
In 2003, the European Committee for Standardization/Information Society Standardization System (CEN/ISSS) DRM Report was published.
In 2004, the Consultation process of the European Commission, and the DG Internal Market, on the Communication COM(2004)261 by the European Commission on "Management of Copyright and Related Rights" closed.
In 2005, DRM Workshops of Directorate-General for Information Society and Media (European Commission), and the work of the High Level Group on DRM were held.
In 2005, Sony BMG installed DRM software on users' computers without clearly notifying the user or requiring confirmation. Among other things, the software included a rootkit, which created a security vulnerability. When the nature of the software was made public much later, Sony BMG initially minimized the significance of the vulnerabilities, but eventually recalled millions of CDs, and made several attempts to patch the software to remove the rootkit. Class action lawsuits were filed, which were ultimately settled by agreements to provide affected consumers with a cash payout or album downloads free of DRM.
Microsoft's media player Zune released in 2006 did not support content that used Microsoft's PlaysForSure DRM scheme.
Windows Media DRM, reads instructions from media files in a rights management language that states what the user may do with the media. Later versions of Windows Media DRM implemented music subscription services that make downloaded files unplayable after subscriptions are cancelled, along with the ability for a regional lockout. Tools like FairUse4WM strip Windows Media of DRM restrictions.
The Gowers Review of Intellectual Property by the British Government from Andrew Gowers was published in 2006 with recommendations regarding copyright terms, exceptions, orphaned works, and copyright enforcement.
DVB (DVB-CPCM) is an updated variant of the broadcast flag. The technical specification was submitted to European governments in March 2007. As with much DRM, the CPCM system is intended to control use of copyrighted material by the end-user, at the direction of the copyright holder. According to Ren Bucholz of the Electronic Frontier Foundation (EFF), "You won't even know ahead of time whether and how you will be able to record and make use of particular programs or devices". The normative sections were approved for publication by the DVB Steering Board, and formalized by ETSI as a formal European Standard (TS 102 825-X) where X refers to the Part number. Nobody has yet stepped forward to provide a Compliance and Robustness regime for the standard, so it is not presently possible to fully implement a system, as no supplier of device certificates has emerged.
In December 2006, the industrial-grade Advanced Access Content System (AACS) for HD DVD and Blu-ray Discs, a process key was published by hackers, which enabled unrestricted access to AACS-protected content.
In January 2007, EMI stopped publishing audio CDs with DRM, stating that "the costs of DRM do not measure up to the results." In March, Musicload.de, one of Europe's largest internet music retailers, announced their position strongly against DRM. In an open letter, Musicload stated that three out of every four calls to their customer support phone service are as a result of consumer frustration with DRM.
Apple Inc. made music DRM-free after April 2007 and labeled all music as "DRM-Free" after 2008. Other works sold on iTunes such as apps, audiobooks, movies, and TV shows are protected by DRM.
A notable DRM failure happened in November 2007, when videos purchased from Major League Baseball prior to 2006 became unplayable due to a change to the servers that validate the licenses.
In 2007, the European Parliament supported the EU's direction on copyright protection.
Asus released a soundcard which features a function called "Analog Loopback Transformation" to bypass the restrictions of DRM. This feature allows the user to record DRM-restricted audio via the soundcard's built-in analog I/O connection.
Digital distributor GOG.com (formerly Good Old Games) specializes in PC video games and has a strict non-DRM policy.
Baen Books and O'Reilly Media, dropped DRM prior to 2012, when Tor Books, a major publisher of science fiction and fantasy books, first sold DRM-free e-books.
The Axmedis project completed in 2008. It was a European Commission Integrated Project of the FP6, has as its main goal automating content production, copy protection, and distribution, to reduce the related costs, and to support DRM at both B2B and B2C areas, harmonizing them.
The INDICARE project was a dialogue on consumer acceptability of DRM solutions in Europe that completed in 2008.
In mid-2008, the Windows version of Mass Effect marked the start of a wave of titles primarily making use of SecuROM for DRM and requiring authentication with a server. The use of the DRM scheme in 2008's Spore led to protests, resulting in searches for an unlicensed version. This backlash against the activation limit led Spore to become the most pirated game in 2008, topping the top 10 list compiled by TorrentFreak. However, Tweakguides concluded that DRM does not appear to increase video game piracy, noting that other games on the list, such as Call of Duty 4 and Assassin's Creed, use DRM without limits or online activation. Additionally, other video games that use DRM, such as BioShock, Crysis Warhead, and Mass Effect, do not appear on the list.
Many mainstream publishers continued to rely on online DRM throughout the later half of 2008 and early 2009, including Electronic Arts, Ubisoft, Valve, and Atari, The Sims 3 being a notable exception in the case of Electronic Arts. Ubisoft broke with the tendency to use online DRM in late 2008, with the release of Prince of Persia as an experiment to "see how truthful people really are" regarding the claim that DRM was inciting people to use illegal copies. Although Ubisoft has not commented on the results of the "experiment", Tweakguides noted that two torrents on Mininova had over 23,000 people downloading the game within 24 hours of its release.
In 2009, Amazon remotely deleted purchased copies of George Orwell's Animal Farm (1945) and Nineteen Eighty-Four (1949) from customers' Amazon Kindles after refunding the purchase price. Commentators described these actions as Orwellian and compared Amazon to Big Brother from Nineteen Eighty-Four. Amazon CEO Jeff Bezos then issued a public apology. FSF wrote that this was an example of the excessive power Amazon has to remotely censor content, and called upon Amazon to drop DRM. Amazon then revealed the reason behind its deletion: the e-books in question were unauthorized reproductions of Orwell's works, which were not within the public domain and that the company that published and sold on Amazon's service had no right to do so.
2010present
Ubisoft formally announced a return to online authentication on 9 February 2010, through its Uplay online game platform, starting with Silent Hunter 5, The Settlers 7, and Assassin's Creed II. Silent Hunter 5 was first reported to have been compromised within 24 hours of release, but users of the cracked version soon found out that only early parts of the game were playable. The Uplay system works by having the installed game on the local PCs incomplete and then continuously downloading parts of the game code from Ubisoft's servers as the game progresses. It was more than a month after the PC release in the first week of April that software was released that could bypass Ubisoft's DRM in Assassin's Creed II. The software did this by emulating a Ubisoft server for the game. Later that month, a real crack was released that was able to remove the connection requirement altogether.
In March 2010, Uplay servers suffered a period of inaccessibility due to a large-scale DDoS attack, causing around 5% of game owners to become locked out of playing their game. The company later credited owners of the affected games with a free download, and there has been no further downtime.
In 2011, comedian Louis C.K. released his concert film Live at the Beacon Theater as an inexpensive (US$5), DRM-free download. The only attempt to deter unlicensed copies was a letter emphasizing the lack of corporate involvement and direct relationship between artist and viewer. The film was a commercial success, turning a profit within 12 hours of its release. The artist suggested that piracy rates were lower than normal as a result, making the release an important case study for the digital marketplace.
In 2012, the EU Court of Justice ruled in favor of reselling copyrighted games.
In 2012, India implemented digital rights management protection.
In 2012, webcomic Diesel Sweeties released a DRM-free PDF e-book. He followed this with a DRM-free iBook specifically for the iPad that generated more than 10,000 downloads in three days. That led Stevens to launch a Kickstarter project – "ebook stravaganza 3000" – to fund the conversion of 3,000 comics, written over 12 years, into a single "humongous" e-book to be released both for free and through the iBookstore; launched 8 February 2012, with the goal of raising $3,000 in 30 days. The "payment optional" DRM-free model in this case was adopted on Stevens' view that "there is a class of webcomics reader who would prefer to read in large chunks and, even better, would be willing to spend a little money on it."
In February 2012, Double Fine asked for crowdfunding for an upcoming video game, Double Fine Adventure, on Kickstarter and offered the game DRM-free for backers. This project exceeded its original goal of $400,000 in 45 days, raising in excess of $2 million. Crowdfunding acted as a pre-order or alternatively as a subscription. After the success of Double Fine Adventure, many games were crowd-funded and many offered a DRM-free version.
Websitessuch as library.nu (shut down by court order on 15 February 2012), BookFi, BookFinder, Library Genesis, and Sci-Huballowed e-book downloading by violating copyright.
As of 2013, other developers, such as Blizzard Entertainment put most of the game logic is on the "side" or taken care of by the servers of the game maker. Blizzard uses this strategy for its game Diablo III and Electronic Arts used this same strategy with their reboot of SimCity, the necessity of which has been questioned.
In 2014, the EU Court of Justice ruled that circumventing DRM on game devices was legal under some circumstances.
In 2014, digital comic distributor Comixology allowed rights holders to provide the option of DRM-free downloads. Publishers that allow this include Dynamite Entertainment, Image Comics, Thrillbent, Top Shelf Productions, and Zenescope Entertainment.
In February 2022, Comixology, which was later under the ownership of Amazon, ended the option of downloading DRM-free downloads on all comics, although any comics previously purchased prior to the date will have the option to download comics without DRM.
Technologies
Verification
Product keys
A product key, typically an alphanumerical string, can represent a license to a particular copy of software. During the installation process or software launch, the user is asked to enter the key; if the key is valid (typically via internal algorithms), the key is accepted, and the user can continue. Product keys can be combined with other DRM practices (such as online "activation"), to prevent cracking the software to run without a product key, or using a keygen to generate acceptable keys.
Activation limits
DRM can limit the number of devices on which a legal user can install content. This restriction typically support 3-5 devices. This affects users who have more devices than the limit. Some allow one device to be replaced with another. Without this software and hardware upgrades may require an additional purchase.
Persistent online DRM
Always-on DRM checks and rechecks authorization while the content is in use by interacting with a server operated by the copyright holder. In some cases, only part of the content is actually installed, while the rest is downloaded dynamically during use.
Encryption
Encryption alters content in a way that means that it cannot be used without first decrypting it. Encryption can ensure that other restriction measures cannot be bypassed by modifying software, so DRM systems typically rely on encryption in addition to other techniques.
Copy restriction
Microsoft PlayReady prevents illicit copying of multimedia and other files.
Restrictions can be applied to electronic books and documents, in order to prevent copying, printing, forwarding, and creating backup copies. This is common for both e-publishers and enterprise Information Rights Management. It typically integrates with content management system software.
While some commentators claim that DRM complicates e-book publishing, it has been used by organizations such as the British Library in its secure electronic delivery service to permit worldwide access to rare documents which, for legal reasons, were previously only available to authorized individuals actually visiting the Library's document centre.
Four main e-book DRM schemes are in common use, from Adobe, Amazon, Apple, and the Marlin Trust Management Organization (MTMO).
Adobe's DRM is applied to EPUBs and PDFs, and can be read by several third-party e-book readers, as well as Adobe Digital Editions (ADE) software. Barnes & Noble uses DRM technology provided by Adobe, applied to EPUBs and the older PDB (Palm OS) format e-books.
Amazon's DRM is an adaption of the original Mobipocket encryption and is applied to Amazon's .azw4, KF8, and Mobipocket format e-books. Topaz format e-books have their own encryption system.
Apple's FairPlay DRM is applied to EPUBs and can be read only by Apple's iBooks app on iOS devices and Mac OS computers.
The Marlin DRM was developed and is maintained by open industry group Marlin Developer Community (MDC) and is licensed by MTMO. (Marlin was founded by Intertrust, Panasonic, Philips, Samsung, and Sony.) Online textbook publisher Kno uses Marlin to protect EPUB books. These books can be read on the Kno App for iOS and Android.
Runtime restrictions
Windows Vista contains a DRM system called Protected Media Path, which contains Protected Video Path (PVP). PVP tries to stop DRM-restricted content from playing while unsigned software is running, in order to prevent the unsigned software from accessing the content. Additionally, PVP can encrypt information during transmission to the monitor or the graphics card, which makes it more difficult to make unauthorized recordings.
Bohemia Interactive have used a form of technology since Operation Flashpoint: Cold War Crisis, wherein if the game copy is suspected of being unauthorized, annoyances like guns losing their accuracy or the players turning into a bird are introduced. Croteam's Serious Sam 3: BFE causes a special invincible foe in the game to appear and constantly attack the player until they are killed.
Regional lockout
Regional lockout (or region coding) prevents the use of a certain product or service, except in a specific region or territory. Lockout may be enforced through physical means, through technological means such as inspecting the user's IP address or using an identifying code, or through unintentional means introduced by devices that support only region-specific technologies (such as video formats, i.e., NTSC and PAL).
Tracking
Watermarks
Digital watermarks can be steganographically embedded within audio or video data. They can be used for recording the copyright owner, the distribution chain or identifying the purchaser. They are not complete DRM mechanisms in their own right, but are used as part of a system for copyright enforcement, such as helping provide evidence for legal purposes, rather than enforcing restrictions.
Some audio/video editing programs may distort, delete, or otherwise interfere with watermarks. Signal/modulator-carrier chromatography may separate watermarks from the recording or detect them as glitches. Additionally, comparison of two separately obtained copies of audio using basic algorithms can reveal watermarks.
Metadata
Sometimes, metadata is included in purchased media which records information such as the purchaser's name, account information, or email address. Also included may be the file's publisher, author, creation date, download date, and various notes. This information is not embedded in the content, as a watermark is. It is kept separate from the content, but within the file or stream.
As an example, metadata is used in media purchased from iTunes for DRM-free as well as DRM-restricted content. This information is included as MPEG standard metadata.
Hardware
US Cable television set-top boxes require a specific piece of hardware to operate. The CableCard standard is used to restrict content to services to which the customer is subscribed. Content has an embedded broadcast flag that the card examines to decide whether the content can be viewed by a specific user.
Implementations
Analog Protection System (Macrovision)
DCS Copy Protection
B-CAS
CableCARD
Broadcast flag
DVB-CPCM
Conditional-access module
Copy Control Information
ISDB#Copy-protection technology
FairPlay
Extended Copy Protection (XCP)
Content Scramble System (CSS)
ARccOS protection
Advanced Access Content System (AACS)
Content Protection for Recordable Media (CPRM)
Digital Transmission Content Protection
High-bandwidth Digital Content Protection (HDCP)
Protected Media Path
Trusted Platform Module#Uses
Intel Management Engine#Design
Cinavia
HTML video Encrypted Media Extensions (HTML EME, often implemented with Widevine)
Denuvo
StarForce
SafeDisc
SecuROM
SafetyNet
Google Play Integrity
In addition, platforms such as Steam may include DRM mechanisms. Most of the mechanisms above are copy protection mechanisms rather than DRM mechanisms per se.
Laws
The World Intellectual Property Organization supports the World Intellectual Property Organization Copyright Treaty (WCT) which requires nations to enact laws against DRM circumvention. The WIPO Internet Treaties do not mandate criminal sanctions, merely requiring "effective legal remedies".
Australia
Australia prohibits circumvention of "access control technical protection measures" in Section 116 of the Copyright Act. The law currently imposes penalties for circumvention of such measures as well as the manufacturing and distribution of tools to enable it.
DRM may be legally circumvented under a few distinct circumstances which are named as exceptions in the law:
permission of the rightsholder
enabling interoperability with copyrighted software
encryption research
security testing
disabling access to private information (circumvention only)
national security or law enforcement
library acquisition decisions (circumvention only)
acts prescribed by regulation (circumvention only)
A person circumventing the access control bears the burden of proof that one of these exceptions apply.
Penalties for violation of the anti-circumvention laws include an injunction, monetary damages, and destruction of enabling devices.
China
China's copyright law was revised in 2001 and included a prohibition on "intentionally circumventing or destroying the technological measures taken by a right holder for protecting the copyright or copyright-related rights in his work, sound recording or video recording, without the permission of the copyright owner, or the owner of the copyright-related rights". However, the Chinese government had faced backlash from Nintendo over the heavy burden
on law enforcement action against circumvention devices, stating that the police only view game copiers as infringing Nintendo's trademark, not as infringing copyright. In response, Nintendo obtained copyright registration for its software in 2013 to make it easier to make law enforcement against game copiers and other circumvention devices.
European Union
The EU operates under its Information Society Directive, its WIPO implementation. The European Parliament then directed member states to outlaw violation of international copyright for commercial purposes. Punishments range from fines to imprisonment. It excluded patent rights and copying for personal, non-commercial purposes. Copyrighted games can be resold. Circumventing DRM on game devices is legal under some circumstances; protections cover only technological measures the interfere with prohibited actions.
India
India acceded to the WIPO Copyright Treaty and the WIPO Performances and Phonograms Treaty on July 4, 2018, after a 2012 amendment to the Copyright Act criminalized the circumvention of technical protections. Fair use is not explicitly addressed, but the anti-circumvention provisions do not prohibit circumventing for non-infringing purposes.
Israel
Israel is not a signatory to the WIPO Copyright Treaty. Israeli law does not expressly prohibit the circumvention of technological protection measures.
Japan
Japan outlawed circumvention of technological protection measures on June 23, 1999 through an amendment of its 1970 copyright law. The private copying exception does not apply if it has become available due to circumvention of TPMs, and circumvention of a TPM is deemed as copyright infringement. However, circumvention is allowed for research purposes or if it otherwise does not harm the rightsholder's interests.
Pakistan
Pakistan is not a signatory to the WIPO Copyright Treaty or the WIPO Performances and Phonograms Treaty. Pakistani law does not criminalize the circumvention of technological protection measures.
As of January 2022, Pakistan's Intellectual Property Office intended to accede to the WIPO Copyright Treaty and WIPO Performances and Phonograms Treaty. However, there has been no major progress for Pakistan to accede to the treaties, and the timeline of the enactments of amendments to the Copyright Ordinance is unclear. As of February 2023, Pakistan's Intellectual Property Office was currently finalizing draft amendments to its Copyright Ordinance.
United States
US protections are governed by the Digital Millennium Copyright Act (DMCA). It criminalizes the production and dissemination of technology that lets users circumvent copy-restrictions. Reverse engineering is expressly permitted, providing a safe harbor where circumvention is necessary to interoperate with other software.
Open-source software that decrypts protected content is not prohibited per se. Decryption done for the purpose of achieving interoperability of open source operating systems with proprietary systems is protected. Dissemination of such software for the purpose of violating or encouraging others to violate copyrights is prohibited.
DMCA has been largely ineffective. Cirumvention software is widely available. However, those who wish to preserve the DRM systems have attempted to use the Act to restrict the distribution and development of such software, as in the case of DeCSS. DMCA contains an exception for research, although the exception is subject to qualifiers that created uncertainty in that community.
Cryptanalytic research may violate the DMCA, although this is unresolved.
Notable lawsuits
DVD Copy Control Association, Inc. v. Bunner
DVD Copy Control Association, Inc. v. Kaleidescape, Inc.
RealNetworks, Inc. v. DVD Copy Control Association, Inc.
Universal v. Reimerdes
Opposition
DRM faces widespread opposition. John Walker and Richard Stallman are notable critics. Stallman has claimed that using the word "rights" is misleading and suggests that the word "restrictions", as in "Digital Restrictions Management", replace it. This terminology has been adopted by other writers and critics.
Other prominent critics include Ross Anderson, who heads a British organization that opposes DRM and similar efforts in the UK and elsewhere, and Cory Doctorow. EFF and organizations such as FreeCulture.org are opposed to DRM. The Foundation for a Free Information Infrastructure criticized DRM's effect as a trade barrier from a free market perspective.
Bruce Schneier argues that digital copy prevention is futile: "What the entertainment industry is trying to do is to use technology to contradict that natural law. They want a practical way to make copying hard enough to save their existing business. But they are doomed to fail." He described trying to make digital files uncopyable as like "trying to make water not wet".
The creators of StarForce stated that "The purpose of copy protection is not making the game uncrackable – it is impossible."
Bill Gates spoke about DRM at 2006 CES, saying that DRM causes problems for legitimate consumers.
The Norwegian consumer rights organization "Forbrukerrådet" complained to Apple in 2007 about the company's use of DRM, accusing it of unlawfully restricting users' access to their music and videos, and of using EULAs that conflict with Norwegian consumer legislation. The complaint was supported by consumers' ombudsmen in Sweden and Denmark, and was reviewed in the EU in 2014. The United States Federal Trade Commission held hearings in March 2009, to review disclosure of DRM limitations to customers' use of media products.
Valve president Gabe Newell stated, "most DRM strategies are just dumb" because they only decrease the value of a game in the consumer's eyes. Newell suggested that the goal should instead be "[creating] greater value for customers through service value". Valve operates Steam, an online store for PC games, as well as a social networking service and a DRM platform.
At the 2012 Game Developers Conference, the CEO of CD Projekt Red, Marcin Iwinski, announced that the company would not use DRM. Iwinski stated of DRM, "It's just over-complicating things... the game... is cracked in two hours." Iwinski added "DRM does not protect your game. If there are examples that it does, then people maybe should consider it, but then there are complications with legit users."
The Association for Computing Machinery and the Institute of Electrical and Electronics Engineers opposed DRM, naming AACS as a technology "most likely to fail" in an issue of IEEE Spectrum.
Public licenses
The GNU General Public License version 3, as released by the Free Software Foundation, has a provision that "strips" DRM of its legal value, so people can break the DRM on GPL software without breaking laws such as the DMCA. In May 2006, FSF launched a "Defective by Design" campaign against DRM.
Creative Commons provides licensing options that encourage creators to work without the use of DRM. Creative Commons licenses have anti-DRM clauses, making the use of DRM by a licensee a breach of the licenses' Baseline Rights.
DRM-free works
Many publishers and artists label their works "DRM-free". Major companies that have done so include Apple, GOG.com, Tor Books and Vimeo on Demand. Comixology once had DRM-free works available for sale until 2022 when its parent company, Amazon, removed the option to buy DRM-free works as part of their migration to Amazon's website, although previous purchases remained DRM-free.
Shortcomings
Availability
Many DRM systems require online authentication. Whenever the server goes down, or a territory experiences an Internet outage, it locks out people from registering or using the material. This is especially true for products that require a persistent online connection, where, for example, a successful DDoS attack on the server essentially makes the material unusable.
Usability
Compact discs (CDs) with DRM schemes are not standards-compliant, and are labeled CD-ROMs. CD-ROMs cannot be played on all CD players or personal computers.
Performance
Certain DRM systems have been associated with reduced performance: some games implementing Denuvo Anti-Tamper performed better without DRM. However, in March 2018, PC Gamer tested Final Fantasy XV for the performance effects of Denuvo, which was found to cause no negative gameplay impact despite a little increase in loading time.
Robustness
DRM copy-prevention schemes can never be wholly secure since the logic needed to decrypt the content is present either in software or hardware and implicitly can be hacked. An attacker can extract this information, decrypt and copy the content, bypassing the DRM.
Satellite and cable systems distribute their content widely and rely on hardware DRM systems. Such systems can be hacked by reverse engineering the protection scheme.
Analog hole
Audio and visual material (excluding interactive materials, e.g., video games) are subject to the analog hole, namely that in order to view the material, the digital signal must be turned into an analog signal. Post-conversion, the material can be then be copied and reconverted to a digital format.
The analog hole cannot be filled without externally imposed restrictions, such as legal regulations, because the vulnerability is inherent to all analog presentation. The conversion from digital to analog and back reduces recording quality. The HDCP attempt to plug the analog hole was largely ineffective.
Consumer rights
Ownership restrictions
DRM opponents argue that it violates private property rights and restricts a range of normal and legal user activities. A DRM component such as that found on a digital audio player restricts how it acts with regard to certain content, overriding user's wishes (for example, preventing the user from copying a copyrighted song to CD as part of a compilation). Doctorow described this as "the right to make up your own copyright laws".
Windows Vista disabled or degraded content play that used a Protected Media Path. DRM restricts the right to make personal copies, provisions lend copies to friends, provisions for service discontinuance, hardware agnosticism, software and operating system agnosticism, lending library use, customer protections against contract amendments by the publisher, and whether content can pass to the owner's heirs.
Obsolescence
When standards and formats change, DRM-restricted content may become obsolete.
When a company undergoes business changes or bankruptcy, its previous services may become unavailable. Examples include MSN Music, Yahoo! Music Store, Adobe Content Server 3 for Adobe PDF, and Acetrax Video on Demand.
Piracy
DRM laws are widely flouted: according to Australia Official Music Chart Survey, copyright infringements from all causes are practised by millions of people. According to the EFF, "in an effort to attract customers, these music services try to obscure the restrictions they impose on you with clever marketing."
Economic implication
Trade-offs between control and sales
Jeff Raikes, ex-president of the Microsoft Business Division, stated: "If they're going to pirate somebody, we want it to be us rather than somebody else". An analogous argument was made in an early paper by Kathleen Conner and Richard Rummelt. A subsequent study of digital rights management for e-books by Gal Oestreicher-Singer and Arun Sundararajan showed that relaxing some forms of DRM can be beneficial to rights holders because the losses from piracy are outweighed by the increase in value to legal buyers. Even if DRM were unbreakable, pirates still might not be willing to purchase, so sales might not increase.
Piracy can be beneficial to some content providers by increase consumer awareness, spreading and popularizing content. This can also increase revenues via other media, such as live performances.
Mathematical models suggest that DRM schemes can fail to do their job on multiple levels. The biggest failure is that the burden that DRM poses on a legitimate customer reduces the customer's willingness to buy. An ideal DRM would not inconvenience legal buyers. The mathematical models are strictly applicable to the music industry.
Alternatives
Several business models offer DRM alternatives.
Subscription
Streaming services have created profitable business models by signing users to monthly subscriptions in return for access to the service's library. This model has worked for music (such as Spotify, Apple Music, etc.) and video (such as Netflix, Disney+, Hulu, etc.).
"Easy and cheap"
Accessing a pirated copy can be illegal and inconvenient. Businesses that charge acceptable fees for doing so tend to attract customers. A business model that dissuades illegal file sharing is to make legal content downloading easy and cheap. Pirate websites often host malware which attaches itself to the files served. If content is provided on legitimate sites and is reasonably priced, consumers are more likely to purchase media legally.
Crowdfunding or pre-order
Crowdfunding has been used as a publishing model for digital content.
Promotion for traditional products
Many artists give away individual tracks to create awareness for a subsequent album.
Artistic Freedom Voucher
The Artistic Freedom Voucher (AFV) introduced by Dean Baker is a way for consumers to support "creative and artistic work". In this system, each consumer receives a refundable tax credit of $100 to give to any artist of creative work. To restrict fraud, the artists must register with the government. The voucher prohibits any artist that receives the benefits from copyrighting their material for a certain length of time. Consumers would be allowed to obtain music for a certain amount of time easily and the consumer would decide which artists receive the $100. The money can either be given to one artist or to many, and this distribution is up to the consumer.
| Technology | Computer security | null |
18940621 | https://en.wikipedia.org/wiki/Ammunition | Ammunition | Ammunition, also known as ammo, is the material fired, scattered, dropped, or detonated from any weapon or weapon system. The term Ammunition includes both expendable weapons (e.g., bombs, missiles, grenades, land mines), and the component parts of other weapons that create the effect on a target (e.g., bullets and warheads).
The purpose of ammunition is to project a force against a selected target to have an effect (usually, but not always, lethal). An example of ammunition is the firearm cartridge, which includes all components required to deliver the weapon effect in a single package. Until the 20th century, black powder was the most common propellant used but has now been replaced in nearly all cases by modern compounds.
Ammunition comes in a great range of sizes and types and is often designed to work only in specific weapons systems. However, there are internationally recognized standards for certain ammunition types (e.g., 5.56×45mm NATO) that enable their use across different weapons and by different users. There are also specific types of ammunition that are designed to have a specialized effect on a target, such as armor-piercing shells and tracer ammunition, used only in certain circumstances. Ammunition is commonly labeled or colored in a specific manner to assist in its identification and to prevent the wrong ammunition types from being used accidentally or inappropriately.
Glossary
A round is a single cartridge containing a projectile, propellant, primer and casing.
A shell is a form of ammunition that is fired by a large caliber cannon or artillery piece. Before the mid-19th century, these shells were usually made of solid materials and relied on kinetic energy to have an effect. However, since that time, they are more often filled with high explosives (see artillery).
A shot is a single release of a weapons system. This may involve firing just one round or piece of ammunition (e.g., from a semi-automatic firearm), but can also refer to ammunition types that release a large number of projectiles at the same time (e.g., cluster munitions or shotgun shells).
A dud is loaded ammunition that fails to function as intended, typically failing to detonate on landing. However, it can also refer to ammunition that fails to fire inside the weapon, known as a misfire, or when the ammunition only partially functions, known as a hang fire. Dud ammunition, which is classified as an unexploded ordnance (UXO), is regarded as highly dangerous. In former conflict zones, it is common for dud ammunition to remain buried in the ground for many years. Large quantities of ammunition from World War I continue to be regularly found in fields throughout France and Belgium and occasionally still claim lives. Although classified as a UXO, landmines that have been left behind after conflict are not considered duds as they have not failed to work and may still be fully functioning.
A bomb or, more specifically, a guided or unguided bomb (also called an aircraft bomb or aerial bomb), is typically an airdropped, unpowered explosive weapon. Mines and the warheads used in guided missiles and rockets are also referred to as bomb-type ammunition.
Etymology
The term ammunition can be traced back to the mid-17th century. The word comes from the French la munition, for the material used for war. Ammunition and munition are often used interchangeably, although munition now usually refers to the actual weapons system with the ammunition required to operate it. In some languages other than English ammunition is still referred to as munition, such as: Dutch ("munitie"), French ("munitions"), German ("Munition"), Italian ("munizione") and Portuguese ("munição").
Design
Ammunition design has evolved throughout history as different weapons have been developed and different effects required. Historically, ammunition was of relatively simple design and build (e.g., sling-shot, stones hurled by catapults), but as weapon designs developed (e.g., rifling) and became more refined, the need for more specialized ammunition increased. Modern ammunition can vary significantly in quality but is usually manufactured to very high standards.
For example, ammunition for hunting can be designed to expand inside a target, maximizing the damage inflicted by one round. Anti-personnel shells are designed to fragment into many pieces and can affect a large area. Armor-piercing rounds are specially hardened to penetrate armor, while smoke ammunition covers an area with a fog that screens people from view. More generic ammunition (e.g., 5.56×45mm NATO) can often be altered slightly to give it a more specific effect (e.g., tracer, incendiary), whilst larger explosive rounds can be altered by using different fuzes.
Components
The components of ammunition intended for rifles and munitions may be divided into these categories:
fuze or primer
explosive materials and propellants
projectiles of all kinds
cartridge casing
Fuzes
The term fuze refers to the detonator of an explosive round or shell. The spelling is different in British English and American English (fuse/fuze respectively) and they are unrelated to a fuse (electrical). A fuse was earlier used to ignite the propellant (e.g., such as on a firework) until the advent of more reliable systems such as the primer or igniter that is used in most modern ammunition.
The fuze of a weapon can be used to alter how the ammunition works. For example, a common artillery shell fuze can be set to "point detonation" (detonation when it hits a target), delay (detonate after it has hit and penetrated a target), time-delay (explode a specified time after firing or impact) and proximity (explode above or next to a target without hitting it, such as for airburst effects or anti-aircraft shells). These allow a single ammunition type to be altered to suit the situation it is required for. There are many designs of a fuze, ranging from simple mechanical to complex radar and barometric systems.
Fuzes are usually armed by the acceleration force of firing the projectile, and usually arm several meters after clearing the bore of the weapon. This helps to ensure the ammunition is safer to handle when loading into the weapon and reduces the chance of the detonator firing before the ammunition has cleared the weapon.
Propellant or explosive
The propellant is the component of ammunition that is activated inside the weapon and provides the kinetic energy required to move the projectile from the weapon to the target. Before the use of gunpowder, this energy would have been produced mechanically by the weapons system (e.g., a catapult or crossbow); in modern times, it is usually a form of chemical energy that rapidly burns to create kinetic force, and an appropriate amount of chemical propellant is packaged with each round of ammunition. In recent years, compressed gas, magnetic energy and electrical energy have been used as propellants.
Until the 20th-century, gunpowder was the most common propellant in ammunition. However, it has since been replaced by a wide range of fast-burning compounds that are more reliable and efficient.
The propellant charge is distinct from the projectile charge which is activated by the fuze, which causes the ammunition effect (e.g., the exploding of an artillery round).
Cartridge case or container
The cartridge is the container that holds the projectile and propellant. Not all ammunition types have a cartridge case. In its place, a wide range of materials can be used to contain the explosives and parts. With some large weapons, the ammunition components are stored separately until loaded into the weapon system for firing. With small arms, caseless ammunition can reduce the weight and cost of ammunition, and simplify the firing process for increased firing rate, but the maturing technology has functionality issues.
Projectile
The projectile is the part of the ammunition that leaves the weapon and has the effect on the target. This effect is usually either kinetic (e.g., as with a standard bullet) or through the delivery of explosives.
Storage
An ammunition dump is a military facility for the storage of live ammunition and explosives that will be distributed and used at a later date. Such a storage facility is extremely hazardous, with the potential for accidents when unloading, packing, and transferring the ammunition. In the event of a fire or explosion, the site and its surrounding area is immediately evacuated and the stored ammunition is left to detonate itself completely with limited attempts at firefighting from a safe distance. In large facilities, there may be a flooding system to automatically extinguish a fire or prevent an explosion. Typically, an ammunition dump will have a large buffer zone surrounding it, to avoid casualties in the event of an accident. There will also be perimeter security measures in place to prevent access by unauthorized personnel and to guard against the potential threat from enemy forces.
A magazine is a place where a quantity of ammunition or other explosive material is stored temporarily prior to being used. The term may be used for a facility where large quantities of ammunition are stored, although this would normally be referred to as an ammunition dump. Magazines are typically located in the field for quick access when engaging the enemy. The ammunition storage area on a warship is referred to as the "ship's magazine". On a smaller scale, magazine is also the name given to the ammunition storage and feeding device of a repeating firearm.
Gunpowder must be stored in a dry place (stable room temperature) to keep it usable, as long as for 10 years. It is also recommended to avoid hot places, because friction or heat might ignite a spark and cause an explosion.
Common types
Small arms
The standard weapon of a modern soldier is an assault rifle, which, like other small arms, uses cartridge ammunition in a size specific to the weapon. Ammunition is carried on the person in box magazines specific to the weapon, ammunition boxes, pouches or bandoliers. The amount of ammunition carried is dependent on the strength of the soldier, the expected action required, and the ability of ammunition to move forward through the logistical chain to replenish the supply. A soldier may also carry a smaller amount of specialized ammunition for heavier weapons such as machine guns and mortars, spreading the burden for squad weapons over many people. Too little ammunition poses a threat to the mission, while too much limits the soldier's mobility also being a threat to the mission.
Shells
A shell is a payload-carrying projectile which, as opposed to a shot, contains explosives or other fillings, in use since the 19th century.
Artillery
Artillery shells are ammunition that is designed to be fired from artillery which has an effect over long distances, usually indirectly (i.e., out of sight of the target). There are many different types of artillery ammunition, but they are usually high-explosive and designed to shatter into fragments on impact to maximize damage. The fuze used on an artillery shell can alter how it explodes or behaves so it has a more specialized effect. Common types of artillery ammunition include high explosive, smoke, illumination, and practice rounds. Some artillery rounds are designed as cluster munitions. Artillery ammunition will almost always include a projectile (the only exception being demonstration or blank rounds), fuze and propellant of some form. When a cartridge case is not used, there will be some other method of containing the propellant bags, usually a breech-loading weapon; see Breechloader.
Tank
Tank ammunition was developed in WWI as tanks first appeared on the battlefield. However, as tank-on-tank warfare developed (including the development of anti-tank warfare artillery), more specialized forms of ammunition were developed such as high-explosive anti-tank (HEAT) warheads and armour-piercing discarding sabot (APDS), including armour-piercing fin-stabilized discarding sabot (APFSDS) rounds. The development of shaped charges has had a significant impact on anti-tank ammunition design, now common in both tank-fired ammunition and in anti-tank missiles, including anti-tank guided missiles.
Naval
Naval weapons were originally the same as many land-based weapons, but the ammunition was designed for specific use, such as a solid shot designed to hole an enemy ship and chain-shot to cut rigging and sails. Modern naval engagements have occurred over far longer distances than historic battles, so as ship armor has increased in strength and thickness, the ammunition to defeat it has also changed. Naval ammunition is now designed to reach very high velocities (to improve its armor-piercing abilities) and may have specialized fuzes to defeat specific types of vessels. However, due to the extended ranges at which modern naval combat may occur, guided missiles have largely supplanted guns and shells.
Aircraft and anti-aircraft
Logistics
With every successive improvement in military arms, a corresponding modification has occurred in the method of supplying ammunition in the quantity required. As soon as projectiles were required (such as javelins and arrows), there needed to be a method of replenishment. When non-specialized, interchangeable or recoverable ammunition was used (e.g., arrows), it was possible to pick up spent arrows (both friendly and enemy) and reuse them. However, with the advent of explosive or non-recoverable ammunition, this was no longer possible and new supplies of ammunition would be needed.
The weight of ammunition required, particularly for artillery shells, can be considerable, causing a need for extra time to replenish supplies. In modern times, there has been an increase in the standardization of many ammunition types between allies (e.g., the NATO Standardization Agreement) that has allowed for shared ammunition types (e.g., 5.56×45mm NATO).
Environmental problems
lead-based ammunition production is the second-largest annual use of lead in the US, accounting for over 60,000 metric tons consumed in 2012. Lead bullets that miss their target or remain in a carcass or body that was never retrieved can enter environmental systems and become toxic to wildlife. The US military has experimented with replacing lead with copper as a slug in their green bullets which reduces the dangers posed by lead in the environment as a result of artillery. Since 2010, this has eliminated over 2000 tons of lead in waste streams.
Hunters are also encouraged to use monolithic bullets, which exclude any lead content.
Unexploded ordnance
Unexploded ammunition can remain active for a very long time and poses a significant threat to both humans and the environment.
| Technology | Ammunition | null |
18943937 | https://en.wikipedia.org/wiki/Emulator | Emulator | In computing, an emulator is hardware or software that enables one computer system (called the host) to behave like another computer system (called the guest). An emulator typically enables the host system to run software or use peripheral devices designed for the guest system.
Emulation refers to the ability of a computer program in an electronic device to emulate (or imitate) another program or device.
Many printers, for example, are designed to emulate HP LaserJet printers because so much software is written for HP printers. If a non-HP printer emulates an HP printer, any software written for a real HP printer will also run in the non-HP printer emulation and produce equivalent printing. Since at least the 1990s, many video game enthusiasts and hobbyists have used emulators to play classic arcade games from the 1980s using the games' original 1980s machine code and data, which is interpreted by a current-era system, and to emulate old video game consoles (see video game console emulator).
A hardware emulator is an emulator which takes the form of a hardware device. Examples include the DOS-compatible card installed in some 1990s-era Macintosh computers, such as the Centris 610 or Performa 630, that allowed them to run personal computer (PC) software programs and field-programmable gate array-based hardware emulators. The Church–Turing thesis implies that theoretically, any operating environment can be emulated within any other environment, assuming memory limitations are ignored. However, in practice, it can be quite difficult, particularly when the exact behavior of the system to be emulated is not documented and has to be deduced through reverse engineering. It also says nothing about timing constraints; if the emulator does not perform as quickly as it did using the original hardware, the software inside the emulation may run much more slowly (possibly triggering timer interrupts that alter behavior).
Types
Most emulators just emulate a hardware architecture—if operating system firmware or software is required for the desired software, it must be provided as well (and may itself be emulated). Both the OS and the software will then be interpreted by the emulator, rather than being run by native hardware. Apart from this interpreter for the emulated binary machine's language, some other hardware (such as input or output devices) must be provided in virtual form as well; for example, if writing to a specific memory location should influence what is displayed on the screen, then this would need to be emulated. While emulation could, if taken to the extreme, go down to the atomic level, basing its output on a simulation of the actual circuitry from a virtual power source, this would be a highly unusual solution. Emulators typically stop at a simulation of the documented hardware specifications and digital logic. Sufficient emulation of some hardware platforms requires extreme accuracy, down to the level of individual clock cycles, undocumented features, unpredictable analog elements, and implementation bugs. This is particularly the case with classic home computers such as the Commodore 64, whose software often depends on highly sophisticated low-level programming tricks invented by game programmers and the "demoscene".
In contrast, some other platforms have had very little use of direct hardware addressing, such as an emulator for the PlayStation 4. In these cases, a simple compatibility layer may suffice. This translates system calls for the foreign system into system calls for the host system e.g., the Linux compatibility layer used on *BSD to run closed source Linux native software on FreeBSD and NetBSD. For example, while the Nintendo 64 graphic processor was fully programmable, most games used one of a few pre-made programs, which were mostly self-contained and communicated with the game via FIFO; therefore, many emulators do not emulate the graphic processor at all, but simply interpret the commands received from the CPU as the original program would. Developers of software for embedded systems or video game consoles often design their software on especially accurate emulators called simulators before trying it on the real hardware. This is so that software can be produced and tested before the final hardware exists in large quantities, so that it can be tested without taking the time to copy the program to be debugged at a low level and without introducing the side effects of a debugger. In many cases, the simulator is actually produced by the company providing the hardware, which theoretically increases its accuracy. Math co-processor emulators allow programs compiled with math instructions to run on machines that do not have the co-processor installed, but the extra work done by the CPU may slow the system down. If a math coprocessor is not installed or present on the CPU, when the CPU executes any co-processor instruction it will make a determined interrupt (coprocessor not available), calling the math emulator routines. When the instruction is successfully emulated, the program continues executing.
Logic simulators
Logic simulation is the use of a computer program to simulate the operation of a digital circuit such as a processor. This is done after a digital circuit has been designed in logic equations, but before the circuit is fabricated in hardware.
Functional emulators
Functional simulation is the use of a computer program to simulate the execution of a second computer program written in symbolic assembly language or compiler language, rather than in binary machine code. By using a functional simulator, programmers can execute and trace selected sections of source code to search for programming errors (bugs), without generating binary code. This is distinct from simulating execution of binary code, which is software emulation. The first functional simulator was written by Autonetics about 1960 for testing assembly language programs for later execution in military computer D-17B. This made it possible for flight programs to be written, executed, and tested before D-17B computer hardware had been built. Autonetics also programmed a functional simulator for testing flight programs for later execution in the military computer D-37C.
Video game console emulators
Video game console emulators are programs that allow a personal computer or video game console to emulate another video game console. They are most often used to play older 1980s to 2000s-era video games on modern personal computers and more contemporary video game consoles. They are also used to translate games into other languages, to modify existing games, and in the development process of "home brew" DIY demos and in the creation of new games for older systems. The Internet has helped in the spread of console emulators, as most - if not all - would be unavailable for sale in retail outlets. Examples of console emulators that have been released in the last few decades are: RPCS3, Dolphin, Cemu, PCSX2, PPSSPP, ZSNES, Citra, ePSXe, Project64, Visual Boy Advance, Nestopia, and Yuzu.
Due to their popularity, emulators have been impersonated by malware. Most of these emulators are for video game consoles like the Xbox 360, Xbox One, Nintendo 3DS, etc. Generally such emulators make currently impossible claims such as being able to run Xbox One and Xbox 360 games in a single program.
Legal issues
As computers and global computer networks continued to advance and emulator developers grew more skilled in their work, the length of time between the commercial release of a console and its successful emulation began to shrink. Fifth generation consoles such as Nintendo 64, PlayStation and sixth generation handhelds, such as the Game Boy Advance, saw significant progress toward emulation during their production. This led to an effort by console manufacturers to stop unofficial emulation, but consistent failures such as Sega v. Accolade 977 F.2d 1510 (9th Cir. 1992), Sony Computer Entertainment, Inc. v. Connectix Corporation 203 F.3d 596 (2000), and Sony Computer Entertainment America v. Bleem 214 F.3d 1022 (2000), have had the opposite effect. According to all legal precedents, emulation is legal within the United States. However, unauthorized distribution of copyrighted code remains illegal, according to both country-specific copyright and international copyright law under the Berne Convention. Under United States law, obtaining a dumped copy of the original machine's BIOS is legal under the ruling Lewis Galoob Toys, Inc. v. Nintendo of America, Inc., 964 F.2d 965 (9th Cir. 1992) as fair use as long as the user obtained a legally purchased copy of the machine. To mitigate this however, several emulators for platforms such as Game Boy Advance are capable of running without a BIOS file, using high-level emulation to simulate BIOS subroutines at a slight cost in emulation accuracy.
Terminal
Terminal emulators are software programs that provide modern computers and devices interactive access to applications running on mainframe computer operating systems or other host systems such as HP-UX or OpenVMS. Terminals such as the IBM 3270 or VT100 and many others are no longer produced as physical devices. Instead, software running on modern operating systems simulates a "dumb" terminal and is able to render the graphical and text elements of the host application, send keystrokes and process commands using the appropriate terminal protocol. Some terminal emulation applications include Attachmate Reflection, IBM Personal Communications, and Micro Focus Rumba.
Other types
Other types of emulators include:
Hardware emulator: the process of imitating the behavior of one or more pieces of hardware (typically a system under design) with another piece of hardware, typically a special purpose emulation system
In-circuit emulator: the use of a hardware device to debug the software of an embedded system
Floating-point emulator: Some floating-point hardware only supports the simplest operations: addition, subtraction, and multiplication. In systems without any floating-point hardware, the CPU emulates it using a series of simpler fixed-point arithmetic operations that run on the integer arithmetic logic unit.
Instruction set simulator in a high-level programming language: Mimics the behavior of a mainframe or microprocessor by "reading" instructions and maintaining internal variables which represent the processor's registers.
Network emulation: a technique for testing the performance of real applications over a virtual network. This is different from network simulation where virtual models of traffic, network models, channels, and protocols are applied.
Server emulator: Multiplayer video games often rely on an online game server, which may or may not be available for on-premises installation. A server emulator is an unofficial on-premises server that imitates the behavior of the official online server, even though its internal working might be different.
Semulation: the process of controlling an emulation through a simulator
Structure and organization
Typically, an emulator is divided into modules that correspond roughly to the emulated computer's subsystems.
Most often, an emulator will be composed of the following modules:
a CPU emulator or CPU simulator (the two terms are mostly interchangeable in this case), unless the target being emulated has the same CPU architecture as the host, in which case a virtual machine layer may be used instead
a memory subsystem module
various input/output (I/O) device emulators
Buses are often not emulated, either for reasons of performance or simplicity, and virtual peripherals communicate directly with the CPU or the memory subsystem.
Memory subsystem
It is possible for the memory subsystem emulation to be reduced to simply an array of elements each sized like an emulated word; however, this model fails very quickly as soon as any location in the computer's logical memory does not match physical memory. This clearly is the case whenever the emulated hardware allows for advanced memory management (in which case, the MMU logic can be embedded in the memory emulator, made a module of its own, or sometimes integrated into the CPU simulator). Even if the emulated computer does not feature an MMU, though, there are usually other factors that break the equivalence between logical and physical memory: many (if not most) architectures offer memory-mapped I/O; even those that do not often have a block of logical memory mapped to ROM, which means that the memory-array module must be discarded if the read-only nature of ROM is to be emulated. Features such as bank switching or segmentation may also complicate memory emulation. As a result, most emulators implement at least two procedures for writing to and reading from logical memory, and it is these procedures' duty to map every access to the correct location of the correct object.
On a base-limit addressing system where memory from address 0 to address ROMSIZE-1 is read-only memory, while the rest is RAM, something along the line of the following procedures would be typical:
void WriteMemory(word Address, word Value) {
word RealAddress;
RealAddress = Address + BaseRegister;
if ((RealAddress < LimitRegister) &&
(RealAddress > ROMSIZE)) {
Memory[RealAddress] = Value;
} else {
RaiseInterrupt(INT_SEGFAULT);
}
}
word ReadMemory(word Address) {
word RealAddress;
RealAddress=Address+BaseRegister;
if (RealAddress < LimitRegister) {
return Memory[RealAddress];
} else {
RaiseInterrupt(INT_SEGFAULT);
return NULL;
}
}
CPU simulator
The CPU simulator is often the most complicated part of an emulator. Many emulators are written using "pre-packaged" CPU simulators, in order to concentrate on good and efficient emulation of a specific machine. The simplest form of a CPU simulator is an interpreter, which is a computer program that follows the execution flow of the emulated program code and, for every machine code instruction encountered, executes operations on the host processor that are semantically equivalent to the original instructions. This is made possible by assigning a variable to each register and flag of the simulated CPU. The logic of the simulated CPU can then more or less be directly translated into software algorithms, creating a software re-implementation that basically mirrors the original hardware implementation.
The following example illustrates how CPU simulation can be accomplished by an interpreter. In this case, interrupts are checked-for before every instruction executed, though this behavior is rare in real emulators for performance reasons (it is generally faster to use a subroutine to do the work of an interrupt).
void Execute(void) {
if (Interrupt != INT_NONE) {
SuperUser = TRUE;
WriteMemory(++StackPointer, ProgramCounter);
ProgramCounter = InterruptPointer;
}
switch (ReadMemory(ProgramCounter++)) {
/*
* Handling of every valid instruction
* goes here...
*/
default:
Interrupt = INT_ILLEGAL;
}
}
Interpreters are very popular as computer simulators, as they are much simpler to implement than more time-efficient alternative solutions, and their speed is more than adequate for emulating computers of more than roughly a decade ago on modern machines. However, the speed penalty inherent in interpretation can be a problem when emulating computers whose processor speed is on the same order of magnitude as the host machine. Until not many years ago, emulation in such situations was considered completely impractical by many.
What allowed breaking through this restriction were the advances in dynamic recompilation techniques. Simple a priori translation of emulated program code into code runnable on the host architecture is usually impossible because of several reasons:
code may be modified while in RAM, even if it is modified only by the emulated operating system when loading the code (for example from disk)
there may not be a way to reliably distinguish data (which should not be translated) from executable code.
Various forms of dynamic recompilation, including the popular Just In Time compiler (JIT) technique, try to circumvent these problems by waiting until the processor control flow jumps into a location containing untranslated code, and only then ("just in time") translates a block of the code into host code that can be executed.
The translated code is kept in a code cache, and the original code is not lost or affected; this way, even data segments can be (meaninglessly) translated by the recompiler, resulting in no more than a waste of translation time. Speed may not be desirable as some older games were not designed with the speed of faster computers in mind. A game designed for a 30 MHz PC with a level timer of 300 game seconds might only give the player 30 seconds on a 300 MHz PC. Other programs, such as some DOS programs, may not even run on faster computers. Particularly when emulating computers which were "closed-box", in which changes to the core of the system were not typical, software may use techniques that depend on specific characteristics of the computer it ran on (e.g. its CPU's speed) and thus precise control of the speed of emulation is important for such applications to be properly emulated.
Input/output (I/O)
Most emulators do not, as mentioned earlier, emulate the main system bus; each I/O device is thus often treated as a special case, and no consistent interface for virtual peripherals is provided. This can result in a performance advantage, since each I/O module can be tailored to the characteristics of the emulated device; designs based on a standard, unified I/O API can, however, rival such simpler models, if well thought-out, and they have the additional advantage of "automatically" providing a plug-in service through which third-party virtual devices can be used within the emulator. A unified I/O API may not necessarily mirror the structure of the real hardware bus: bus design is limited by several electric constraints and a need for hardware concurrency management that can mostly be ignored in a software implementation.
Even in emulators that treat each device as a special case, there is usually a common basic infrastructure for:
managing interrupts, by means of a procedure that sets flags readable by the CPU simulator whenever an interrupt is raised, allowing the virtual CPU to "poll for (virtual) interrupts"
writing to and reading from physical memory, by means of two procedures similar to the ones dealing with logical memory (although, contrary to the latter, the former can often be left out, and direct references to the memory array be employed instead)
Applications
In preservation
Emulation is one strategy in pursuit of digital preservation and combating obsolescence. Emulation focuses on recreating an original computer environment, which can be time-consuming and difficult to achieve, but valuable because of its ability to maintain a closer connection to the authenticity of the digital object, operating system, or even gaming platform. Emulation addresses the original hardware and software environment of the digital object, and recreates it on a current machine. The emulator allows the user to have access to any kind of application or operating system on a current platform, while the software runs as it did in its original environment. Jeffery Rothenberg, an early proponent of emulation as a digital preservation strategy states, "the ideal approach would provide a single extensible, long-term solution that can be designed once and for all and applied uniformly, automatically, and in organized synchrony (for example, at every refresh cycle) to all types of documents and media". He further states that this should not only apply to out of date systems, but also be upwardly mobile to future unknown systems. Practically speaking, when a certain application is released in a new version, rather than address compatibility issues and migration for every digital object created in the previous version of that application, one could create an emulator for the application, allowing access to all of said digital objects.
In new media art
Because of its primary use of digital formats, new media art relies heavily on emulation as a preservation strategy. Artists such as Cory Arcangel specialize in resurrecting obsolete technologies in their artwork and recognize the importance of a decentralized and deinstitutionalized process for the preservation of digital culture. In many cases, the goal of emulation in new media art is to preserve a digital medium so that it can be saved indefinitely and reproduced without error, so that there is no reliance on hardware that ages and becomes obsolete. The paradox is that the emulation and the emulator have to be made to work on future computers.
In future systems design
Emulation techniques are commonly used during the design and development of new systems. It eases the development process by providing the ability to detect, recreate and repair flaws in the design even before the system is actually built. It is particularly useful in the design of multi-core systems, where concurrency errors can be very difficult to detect and correct without the controlled environment provided by virtual hardware. This also allows the software development to take place before the hardware is ready, thus helping to validate design decisions and give a little more control.
Comparison with simulation
The word "emulator" was coined in 1963 at IBM during development of the NPL (IBM System/360) product line, using a "new combination of software, microcode, and hardware". They discovered that simulation using additional instructions implemented in microcode and hardware, instead of software simulation using only standard instructions, to execute programs written for earlier IBM computers dramatically increased simulation speed. Earlier, IBM provided simulators for, e.g., the 650 on the 705. In addition to simulators, IBM had compatibility features on the 709 and 7090, for which it provided the IBM 709 computer with a program to run legacy programs written for the IBM 704 on the 709 and later on the IBM 7090. This program used the instructions added by the compatibility feature to trap instructions requiring special handling; all other 704 instructions ran the same on a 7090. The compatibility feature on the 1410 only required setting a console toggle switch, not a support program.
In 1963, when microcode was first used to speed up this simulation process, IBM engineers coined the term "emulator" to describe the concept. In the 2000s, it has become common to use the word "emulate" in the context of software. However, before 1980, "emulation" referred only to emulation with a hardware or microcode assist, while "simulation" referred to pure software emulation. For example, a computer specially built for running programs designed for another architecture is an emulator. In contrast, a simulator could be a program which runs on a PC, so that old Atari games can be simulated on it. Purists continue to insist on this distinction, but currently the term "emulation" often means the complete imitation of a machine executing binary code while "simulation" often refers to computer simulation, where a computer program is used to simulate an abstract model. Computer simulation is used in virtually every scientific and engineering domain and Computer Science is no exception, with several projects simulating abstract models of computer systems, such as network simulation, which both practically and semantically differs from network emulation.
Comparison with hardware virtualization
Hardware virtualization is the virtualization of computers as complete hardware platforms, certain logical abstractions of their components, or only the functionality required to run various operating systems. Virtualization hides the physical characteristics of a computing platform from the users, presenting instead an abstract computing platform. At its origins, the software that controlled virtualization was called a "control program", but the terms "hypervisor" or "virtual machine monitor" became preferred over time. Each hypervisor can manage or run multiple virtual machines.
| Technology | Computer software | null |
18947686 | https://en.wikipedia.org/wiki/Caryophyllales | Caryophyllales | Caryophyllales ( ) is a diverse and heterogeneous order of flowering plants that includes the cacti, carnations, amaranths, ice plants, beets, and many carnivorous plants. Many members are succulent, having fleshy stems or leaves. The betalain pigments are unique in plants of this order and occur in all its core families with the exception of Caryophyllaceae and Molluginaceae. Noncore families, such as Nepenthaceae, instead produce anthocyanins.
In its modern definition, the order encompasses a whole new group of families (formerly included in the order Polygonales) that never synthesize betalains, among which several families are carnivorous (like Nepenthaceae and Droseraceae).
According to molecular clock calculations, the lineage that led to Caryophyllales split from other plants about 111 million years ago.
Description
The members of Caryophyllales include about 6% of eudicot species. This order is part of the core eudicots. Currently, the Caryophyllales contains 37 families, 749 genera, and 11,620 species The monophyly of the Caryophyllales has been supported by DNA sequences, cytochrome c sequence data and heritable characters such as anther wall development and vessel-elements with simple perforations.
Circumscription
As with all taxa, the circumscription of Caryophyllales has changed within various classification systems. All systems recognize a core of families with centrospermous ovules and seeds. More recent treatments have expanded the Caryophyllales to include many carnivorous plants.
Systematists were undecided on whether Caryophyllales should be placed within the rosid complex or sister to the asterid clade. The possible connection between sympetalous angiosperms and Caryophyllales was presaged by Bessey, Hutchinson, and others; as Lawrence relates: "The evidence is reasonably conclusive that the Primulaceae and the Caryophyllaceae have fundamentally the same type of gynecia, and as concluded by Douglas (1936)(and essentially Dickson, 1936) '...the vascular pattern and the presence of locules at the base of the ovary point to the fact that the present much reduced flower of the Primulaceae has descended from an ancestor which was characterized by a plurilocular ovary and axial placentation. This primitive flower might well be found in centrospermal stock as Wernham, Bessy, and Hutchinson have suggested.' "
Caryophyllales is separated into two suborders: Caryophyllineae and Polygonineae. These two suborders were formerly (and sometimes still are) recognized as two orders, Polygonales and Caryophyllales.
APG IV
Kewaceae, Macarthuriaceae, Microteaceae, and Petiveriaceae were added in APG IV.
APG III
As circumscribed by the APG III system (2009), this order includes the same families as the APG II system (see below) plus the new families, Limeaceae, Lophiocarpaceae, Montiaceae, Talinaceae, and Anacampserotaceae.
family Achatocarpaceae
family Aizoaceae
family Amaranthaceae
family Anacampserotaceae
family Ancistrocladaceae
family Asteropeiaceae
family Barbeuiaceae
family Basellaceae
family Cactaceae
family Caryophyllaceae
family Didiereaceae
family Dioncophyllaceae
family Droseraceae
family Drosophyllaceae
family Frankeniaceae
family Gisekiaceae
family Halophytaceae
family Kewaceae
family Limeaceae
family Lophiocarpaceae
family Macarthuriaceae
family Microteaceae
family Molluginaceae
family Montiaceae
family Nepenthaceae
family Nyctaginaceae
family Petiveriaceae
family Physenaceae
family Phytolaccaceae
family Plumbaginaceae
family Polygonaceae
family Portulacaceae
family Rhabdodendraceae
family Sarcobataceae
family Simmondsiaceae
family Stegnospermataceae
family Talinaceae
family Tamaricaceae
APG II
As circumscribed by the APG II system (2003), this order includes well-known plants like cacti, carnations, spinach, beet, rhubarb, sundews, venus fly traps, and bougainvillea. Recent molecular and biochemical evidence has resolved additional well-supported clades within the Caryophyllales.
order Caryophyllales
family Achatocarpaceae
family Aizoaceae
family Amaranthaceae
family Anacampserotaceae (added in APG III)
family Ancistrocladaceae
family Asteropeiaceae
family Barbeuiaceae
family Basellaceae
family Cactaceae
family Caryophyllaceae
family Didiereaceae
family Dioncophyllaceae
family Droseraceae
family Drosophyllaceae
family Frankeniaceae
family Gisekiaceae
family Halophytaceae
family Limeaceae (added in APG III)
family Lophiocarpaceae (added in APG III)
family Molluginaceae
family Montiaceae (added in APG III)
family Nepenthaceae
family Nyctaginaceae
family Physenaceae
family Phytolaccaceae
family Plumbaginaceae
family Polygonaceae
family Portulacaceae
family Rhabdodendraceae
family Sarcobataceae
family Simmondsiaceae
family Stegnospermataceae
family Talinaceae (added in APG III)
family Tamaricaceae
APG
This represents a slight change from the APG system, of 1998
order Caryophyllales
family Achatocarpaceae
family Aizoaceae
family Amaranthaceae
family Ancistrocladaceae
family Asteropeiaceae
family Basellaceae
family Cactaceae
family Caryophyllaceae
family Didiereaceae
family Dioncophyllaceae
family Droseraceae
family Drosophyllaceae
family Frankeniaceae
family Molluginaceae
family Nepenthaceae
family Nyctaginaceae
family Physenaceae
family Phytolaccaceae
family Plumbaginaceae
family Polygonaceae
family Portulacaceae
family Rhabdodendraceae
family Sarcobataceae
family Simmondsiaceae
family Stegnospermataceae
family Tamaricaceae
Cronquist
The Cronquist system (1981) also recognised the order, with this circumscription:
order Caryophyllales
family Achatocarpaceae
family Aizoaceae
family Amaranthaceae
family Basellaceae
family Cactaceae
family Caryophyllaceae
family Chenopodiaceae
family Didiereaceae
family Nyctaginaceae
family Phytolaccaceae
family Portulacaceae
family Molluginaceae
The difference with the order as recognized by APG lies in the first place in the concept of "order". The APG favours much larger orders and families, and the order Caryophyllales sensu APG should rather be compared to subclass Caryophyllidae sensu Cronquist.
A part of the difference lies with what families are recognized. The plants in the Stegnospermataceae and Barbeuiaceae were included in Cronquist's Phytolaccaceae. The Chenopodiaceae (still recognized by Cronquist) are included in Amaranthaceae by APG.
New to the order (sensu APG) are the Asteropeiaceae and Physenaceae, each containing a single genus, and two genera from Cronquist's order Nepenthales.
Earlier circumscriptions
Earlier systems, such as the Wettstein system, last edition in 1935, and the Engler system, updated in 1964, had a similar order under the name Centrospermae.
| Biology and health sciences | Caryophyllales | Plants |
18947703 | https://en.wikipedia.org/wiki/Nausea | Nausea | Nausea is a diffuse sensation of unease and discomfort, sometimes perceived as an urge to vomit. It can be a debilitating symptom if prolonged and has been described as placing discomfort on the chest, abdomen, or back of the throat.
Over 30 definitions of nausea were proposed in a 2011 book on the topic.
Nausea is a non-specific symptom, which means that it has many possible causes. Some common causes of nausea are gastroenteritis and other gastrointestinal disorders, food poisoning, motion sickness, dizziness, migraine, fainting, low blood sugar, anxiety, hyperthermia, dehydration and lack of sleep. Nausea is a side effect of many medications including chemotherapy, or morning sickness in early pregnancy. Nausea may also be caused by disgust and depression.
Medications taken to prevent and treat nausea and vomiting are called antiemetics. The most commonly prescribed antiemetics in the US are promethazine, metoclopramide, and the newer ondansetron. The word nausea is from Latin nausea, from Greek – nausia, "ναυτία" – nautia, motion sickness, "feeling sick or queasy".
Causes
Gastrointestinal infections (37%) and food poisoning are the two most common causes of acute nausea and vomiting. Side effects from medications (3%) and pregnancy are also relatively frequent. There are many causes of chronic nausea. Nausea and vomiting remain undiagnosed in 10% of the cases. Aside from morning sickness, there are no sex differences in complaints of nausea. After childhood, doctor consultations decrease steadily with age. Only a fraction of one percent of doctor visits by those over 65 are due to nausea.
Gastrointestinal
Gastrointestinal infection is one of the most common causes of acute nausea and vomiting. Chronic nausea may be the presentation of many gastrointestinal disorders, occasionally as the major symptom, such as gastroesophageal reflux disease, functional dyspepsia, gastritis, biliary reflux, gastroparesis, peptic ulcer, celiac disease, non-celiac gluten sensitivity, Crohn's disease, hepatitis, upper gastrointestinal malignancy, and pancreatic cancer. Uncomplicated Helicobacter pylori infection does not cause chronic nausea.
Food poisoning
Food poisoning usually causes an abrupt onset of nausea and vomiting one to six hours after ingestion of contaminated food and lasts for one to two days. It is due to toxins produced by bacteria in food.
Medications
Many medications can potentially cause nausea. Some of the most frequently associated include cytotoxic chemotherapy regimens for cancer and other diseases, and general anaesthetic agents. An old cure for migraine, ergotamine, is well known to cause devastating nausea in some patients; a person using it for the first time will be prescribed an antiemetic for relief if needed.
Pregnancy
Nausea or "morning sickness" is common during early pregnancy but may occasionally continue into the second and third trimesters. In the first trimester nearly 80 % of women have some degree of nausea. Pregnancy should therefore be considered as a possible cause of nausea in any sexually active woman of child-bearing age. While usually it is mild and self-limiting, severe cases known as hyperemesis gravidarum may require treatment.
Disequilibrium
A number of conditions involving balance such as motion sickness and vertigo can lead to nausea and vomiting.
Gynecologic
Dysmenorrhea can cause nausea.
Psychiatric
Nausea may be caused by depression, anxiety disorders and eating disorders.
Potentially serious
While most causes of nausea are not serious, some serious conditions are associated with nausea. These include pancreatitis, small bowel obstruction, appendicitis, cholecystitis, hepatitis, Addisonian crisis, diabetic ketoacidosis, increased intracranial pressure, spontaneous intracranial hypotension, brain tumors, meningitis, heart attack, rabies, carbon monoxide poisoning and many others.
Comprehensive list
Inside the abdomen
Obstructing disorders
Gastric outlet obstruction
Small bowel obstruction
Colonic obstruction
Superior mesenteric artery syndrome
Enteric infections
Viral infection
Bacterial infection
Inflammatory diseases
Celiac disease
Cholecystitis
Pancreatitis
Appendicitis
Hepatitis
Sensorimotor dysfunction
Gastroparesis
Intestinal pseudo-obstruction
Gastroesophageal reflux disease
Irritable bowel syndrome
Cyclic vomiting syndrome
Other
Non-celiac gluten sensitivity
Biliary colic
Kidney stone
Cirrhosis
Abdominal irradiation
Outside the abdomen
Cardiopulmonary
Cardiomyopathy
Myocardial infarction (heart attack)
Paroxysmal cough
Inner-ear diseases
Motion sickness
Labyrinthitis
Malignancy
Intracerebral disorders
Malignancy
Hemorrhage
Abscess
Hydrocephalus
Meningitis
Encephalitis
Rabies
Psychiatric illnesses
Anorexia and bulimia nervosa
Depression
Drug withdrawal
Other
Post-operative vomiting
Nociception
Altitude sickness
Medications and metabolic disorders
Drugs
Chemotherapy
Antibiotics
Antiarrhythmics
Digoxin
Oral hypoglycemic medications
Oral contraceptives
Norepinephrine reuptake inhibitors
Endocrine/metabolic disease
Pregnancy
Uremia
Ketoacidosis
Thyroid and parathyroid disease
Adrenal insufficiency
Toxins
Liver failure
Alcohol
Pathophysiology
Research on nausea and vomiting has relied on using animal models to mimic the anatomy and neuropharmacologic features of the human body. The physiologic mechanism of nausea is a complex process that has yet to be fully elucidated. There are four general pathways that are activated by specific triggers in the human body that go on to create the sensation of nausea and vomiting.
Central nervous system (CNS): Stimuli can affect areas of the CNS including the cerebral cortex and the limbic system. These areas are activated by elevated intracranial pressure, irritation of the meninges (i.e. blood or infection), and extreme emotional triggers such as anxiety. The supratentorial region is also responsible for the sensation of nausea.
Chemoreceptor trigger zone (CTZ): The CTZ is located in the area postrema in the floor of the fourth ventricle within the brain. This area is outside the blood brain barrier, and is therefore readily exposed to substances circulating through the blood and cerebral spinal fluid. Common triggers of the CTZ include metabolic abnormalities, toxins, and medications. Activation of the CTZ is mediated by dopamine (D2) receptors, serotonin (5HT3) receptors, and neurokinin receptors (NK1).
Vestibular system: This system is activated by disturbances to the vestibular apparatus in the inner ear. These include movements that cause motion sickness and dizziness. This pathway is triggered via histamine (H1) receptors and acetylcholine (ACh) receptors.
Peripheral Pathways: These pathways are triggered via chemoreceptors and mechanoreceptors in the gastrointestinal tract, as well as other organs such as the heart and kidneys. Common activators of these pathways include toxins present in the gastrointestinal lumen and distension of the gastrointestinal lumen from blockage or dysmotility of the bowels. Signals from these pathways travel via multiple neural tracts including the vagus, glossopharyngeal, splanchnic, and sympathetic nerves.
Signals from any of these pathways then travel to the brainstem, activating several structures including the nucleus of the solitary tract, the dorsal motor nucleus of the vagus, and central pattern generator. These structures go on to signal various downstream effects of nausea and vomiting. The body's motor muscle responses involve halting the muscles of the gastrointestinal tract, and in fact causing reversed propulsion of gastric contents towards the mouth while increasing abdominal muscle contraction. Autonomic effects involve increased salivation and the sensation of feeling faint that often occurs with nausea and vomiting.
Pre-nausea pathophysiology
It has been described that alterations in heart rate can occur as well as the release of vasopressin from the posterior pituitary.
Diagnosis
Patient history
Taking a thorough patient history may reveal important clues to the cause of nausea and vomiting. If the patient's symptoms have an acute onset, then drugs, toxins, and infections are likely. In contrast, a long-standing history of nausea will point towards a chronic illness as the culprit. The timing of nausea and vomiting after eating food is an important factor to pay attention to. Symptoms that occur within an hour of eating may indicate an obstruction proximal to the small intestine, such as gastroparesis or pyloric stenosis. An obstruction further down in the intestine or colon will cause delayed vomiting. An infectious cause of nausea and vomiting such as gastroenteritis may present several hours to days after the food was ingested. The contents of the emesis is a valuable clue towards determining the cause. Bits of fecal matter in the emesis indicate obstruction in the distal intestine or the colon. Emesis that is of a bilious nature (greenish in color) localizes the obstruction to a point past the stomach. Emesis of undigested food points to an obstruction prior to the gastric outlet, such as achalasia or Zenker's diverticulum. If patient experiences reduced abdominal pain after vomiting, then obstruction is a likely etiology. However, vomiting does not relieve the pain brought on by pancreatitis or cholecystitis.
Physical exam
It is important to watch out for signs of dehydration, such as orthostatic hypotension and loss of skin turgor. Auscultation of the abdomen can produce several clues to the cause of nausea and vomiting. A high-pitched tinkling sound indicates possible bowel obstruction, while a splashing "succussion" sound is more indicative of gastric outlet obstruction. Eliciting pain on the abdominal exam when pressing on the patient may indicate an inflammatory process. Signs such as papilledema, visual field losses, or focal neurological deficits are red flag signs for elevated intracranial pressure.
Diagnostic testing
When a history and physical exam are not enough to determine the cause of nausea and vomiting, certain diagnostic tests may prove useful. A chemistry panel would be useful for electrolyte and metabolic abnormalities. Liver function tests and lipase would identify pancreaticobiliary diseases. Abdominal X-rays showing air-fluid levels indicate bowel obstruction, while an X-ray showing air-filled bowel loops are more indicative of ileus. More advanced imaging and procedures may be necessary, such as a CT scan, upper endoscopy, colonoscopy, barium enema, or MRI. Abnormal GI motility can be assessed using specific tests like gastric scintigraphy, wireless motility capsules, and small-intestinal manometry.
Treatment
If dehydration is present due to loss of fluids from severe vomiting, rehydration with oral electrolyte solutions is preferred. If this is not effective or possible, intravenous rehydration may be required. Medical care is recommended if: a person cannot keep any liquids down, has symptoms more than 2 days, is weak, has a fever, has stomach pain, vomits more than two times in a day or does not urinate for more than 8 hours.
Medications
Numerous pharmacologic medications are available for the treatment of nausea. There is no medication that is clearly superior to other medications for all cases of nausea. The choice of antiemetic medication may be based on the situation during which the person experiences nausea. For people with motion sickness and vertigo, antihistamines and anticholinergics such as meclizine and scopolamine are particularly effective. Nausea and vomiting associated with migraine headaches respond best to dopamine antagonists such as metoclopramide, prochlorperazine, and chlorpromazine. In cases of gastroenteritis, serotonin antagonists such as ondansetron were found to suppress nausea and vomiting, as well as reduce the need for IV fluid resuscitation. The combination of pyridoxine and doxylamine is the first line treatment for pregnancy-related nausea and vomiting. Dimenhydrinate is an inexpensive and effective over the counter medication for preventing postoperative nausea and vomiting. Other factors to consider when choosing an antiemetic medication include the person's preference, side-effect profile, and cost.
Nabilone is also indicated for this purpose.
Alternative medicine
In certain people, cannabinoids may be effective in reducing chemotherapy associated nausea and vomiting. Several studies have demonstrated the therapeutic effects of cannabinoids for nausea and vomiting in the advanced stages of illnesses such as cancer and AIDS.
In hospital settings topical anti-nausea gels are not indicated because of lack of research backing their efficacy. Topical gels containing lorazepam, diphenhydramine, and haloperidol are sometimes used for nausea but are not equivalent to more established therapies.
Ginger has also been shown to be potentially effective in treating several types of nausea.
Prognosis
The outlook depends on the cause. Most people recover within few hours or a day. While short-term nausea and vomiting are generally harmless, they may sometimes indicate a more serious condition. When associated with prolonged vomiting, it may lead to dehydration or dangerous electrolyte imbalances or both. Repeated intentional vomiting, characteristic of bulimia, can cause stomach acid to wear away at the enamel present on the teeth.
Epidemiology
Nausea and or vomiting is the main complaint in 1.6% of visits to family physicians in Australia. However, only 25% of people with nausea visit their family physician. In Australia, nausea, as opposed to vomiting, occurs most frequently in persons aged 15–24 years, and is less common in other age groups.
| Biology and health sciences | Symptoms and signs | Health |
27471338 | https://en.wikipedia.org/wiki/Object-oriented%20programming | Object-oriented programming | Object-oriented programming (OOP) is a programming paradigm based on the concept of objects, which can contain data and code: data in the form of fields (often known as attributes or properties), and code in the form of procedures (often known as methods). In OOP, computer programs are designed by making them out of objects that interact with one another.
Many of the most widely used programming languages (such as C++, Java, and Python) are multi-paradigm and support object-oriented programming to a greater or lesser degree, typically in combination with imperative programming and declarative programming.
Significant object-oriented languages include Ada, ActionScript, C++, Common Lisp, C#, Dart, Eiffel, Fortran 2003, Haxe, Java, JavaScript, Kotlin, Logo, MATLAB, Objective-C, Object Pascal, Perl, PHP, Python, R, Raku, Ruby, Scala, SIMSCRIPT, Simula, Smalltalk, Swift, Vala and Visual Basic.NET.
History
Terminology invoking "objects" in the modern sense of object-oriented programming made its first appearance at the artificial intelligence group at MIT in the late 1950s and early 1960s. "Object" referred to LISP atoms with identified properties (attributes).
Another early MIT example was Sketchpad created by Ivan Sutherland in 1960–1961; in the glossary of the 1963 technical report based on his dissertation about Sketchpad, Sutherland defined notions of "object" and "instance" (with the class concept covered by "master" or "definition"), albeit specialized to graphical interaction.
Also, in 1968, an MIT ALGOL version, AED-0, established a direct link between data structures ("plexes", in that dialect) and procedures, prefiguring what were later termed "messages", "methods", and "member functions".
Topics such as data abstraction and modular programming were common points of discussion at this time.
Independently of later MIT work such as AED, Simula was developed during the years 1961–1967.
Simula introduced important concepts that are today an essential part of object-oriented programming, such as class and object, inheritance, and dynamic binding.
The object-oriented Simula programming language was used mainly by researchers involved with physical modelling, such as models to study and improve the movement of ships and their content through cargo ports.
Influenced by the work at MIT and the Simula language, in November 1966 Alan Kay began working on ideas that would eventually be incorporated into the Smalltalk programming language. Kay used the term "object-oriented programming" in conversation as early as 1967. Although sometimes called "the father of object-oriented programming", Alan Kay has differentiated his notion of OO from the more conventional abstract data type notion of object, and has implied that the computer science establishment did not adopt his notion.
A 1976 MIT memo co-authored by Barbara Liskov lists Simula 67, CLU, and Alphard as object-oriented languages, but does not mention Smalltalk.
In the 1970s, the first version of the Smalltalk programming language was developed at Xerox PARC by Alan Kay, Dan Ingalls and Adele Goldberg. Smalltalk-72 included a programming environment and was dynamically typed, and at first was interpreted, not compiled.
Smalltalk became noted for its application of object orientation at the language-level and its graphical development environment. Smalltalk went through various versions and interest in the language grew.
While Smalltalk was influenced by the ideas introduced in Simula 67 it was designed to be a fully dynamic system in which classes could be created and modified dynamically.
During the late 1970s and 1980s, object-oriented programming rose to prominence. The Flavors object-oriented Lisp was developed starting 1979, introducing multiple inheritance and mixins. In 1981, Goldberg edited the August issue of Byte Magazine, introducing Smalltalk and object-oriented programming to a wide audience. LOOPS, the object system for Interlisp-D, was influenced by Smalltalk and Flavors, and a paper about it was published in 1982. In 1986, the Association for Computing Machinery organized the first Conference on Object-Oriented Programming, Systems, Languages, and Applications (OOPSLA), which was attended by 1,000 people. Among other developments was the Common Lisp Object System, which integrates functional programming and object-oriented programming and allows extension via a Meta-object protocol. In the 1980s, there were a few attempts to design processor architectures that included hardware support for objects in memory but these were not successful. Examples include the Intel iAPX 432 and the Linn Smart Rekursiv.
In the mid-1980s Objective-C was developed by Brad Cox, who had used Smalltalk at ITT Inc.. Bjarne Stroustrup, who had used Simula for his PhD thesis, created the object-oriented C++. In 1985, Bertrand Meyer also produced the first design of the Eiffel language. Focused on software quality, Eiffel is a purely object-oriented programming language and a notation supporting the entire software lifecycle. Meyer described the Eiffel software development method, based on a small number of key ideas from software engineering and computer science, in Object-Oriented Software Construction. Essential to the quality focus of Eiffel is Meyer's reliability mechanism, design by contract, which is an integral part of both the method and language.
In the early and mid-1990s object-oriented programming developed as the dominant programming paradigm when programming languages supporting the techniques became widely available. These included Visual FoxPro 3.0, C++, and Delphi. Its dominance was further enhanced by the rising popularity of graphical user interfaces, which rely heavily upon object-oriented programming techniques. An example of a closely related dynamic GUI library and OOP language can be found in the Cocoa frameworks on Mac OS X, written in Objective-C, an object-oriented, dynamic messaging extension to C based on Smalltalk. OOP toolkits also enhanced the popularity of event-driven programming (although this concept is not limited to OOP).
At ETH Zürich, Niklaus Wirth and his colleagues investigated the concept of type checking across module boundaries. Modula-2 (1978) included this concept, and their succeeding design, Oberon (1987), included a distinctive approach to object orientation, classes, and such. Inheritance is not obvious in Wirth's design since his nomenclature looks in the opposite direction: It is called type extension and the viewpoint is from the parent down to the inheritor.
Object-oriented features have been added to many previously existing languages, including Ada, BASIC, Fortran, Pascal, and COBOL. Adding these features to languages that were not initially designed for them often led to problems with compatibility and maintainability of code.
More recently, some languages have emerged that are primarily object-oriented, but that are also compatible with procedural methodology. Two such languages are Python and Ruby. Probably the most commercially important recent object-oriented languages are Java, developed by Sun Microsystems, as well as C# and Visual Basic.NET (VB.NET), both designed for Microsoft's .NET platform. Each of these two frameworks shows, in its way, the benefit of using OOP by creating an abstraction from implementation. VB.NET and C# support cross-language inheritance, allowing classes defined in one language to subclass classes defined in the other language.
Features
Object-oriented programming uses objects, but not all of the associated techniques and structures are supported directly in languages that claim to support OOP. The features listed below are common among languages considered to be strongly class- and object-oriented (or multi-paradigm with OOP support), with notable exceptions mentioned. Christopher J. Date stated that critical comparison of OOP to other technologies, relational in particular, is difficult because of lack of an agreed-upon and rigorous definition of OOP.
Shared with non-OOP languages
Variables that can store information formatted in a small number of built-in data types like integers and alphanumeric characters. This may include data structures like strings, lists, and hash tables that are either built-in or result from combining variables using memory pointers.
Procedures – also known as functions, methods, routines, or subroutines – that take input, generate output, and manipulate data. Modern languages include structured programming constructs like loops and conditionals.
Modular programming support provides the ability to group procedures into files and modules for organizational purposes. Modules are namespaced so identifiers in one module will not conflict with a procedure or variable sharing the same name in another file or module.
Objects
An object is a data structure or abstract data type containing fields (state variables containing data) and methods (subroutines or procedures defining the object's behavior in code). Fields may also be known as members, attributes, or properties. Objects are typically stored as contiguous regions of memory. Objects are accessed somewhat like variables with complex internal structures, and in many languages are effectively pointers, serving as actual references to a single instance of said object in memory within a heap or stack.
Objects sometimes correspond to things found in the real world. For example, a graphics program may have objects such as "circle", "square", and "menu". An online shopping system might have objects such as "shopping cart", "customer", and "product". Sometimes objects represent more abstract entities, like an object that represents an open file, or an object that provides the service of translating measurements from U.S. customary to metric.
Objects can contain other objects in their instance variables; this is known as object composition. For example, an object in the Employee class might contain (either directly or through a pointer) an object in the Address class, in addition to its own instance variables like "first_name" and "position". Object composition is used to represent "has-a" relationships: every employee has an address, so every Employee object has access to a place to store an Address object (either directly embedded within itself or at a separate location addressed via a pointer). Date and Darwen have proposed a theoretical foundation that uses OOP as a kind of customizable type system to support RDBMS, but it forbids object pointers.
The OOP paradigm has been criticized for overemphasizing the use of objects for software design and modeling at the expense of other important aspects (computation/algorithms). For example, Rob Pike has said that OOP languages frequently shift the focus from data structures and algorithms to types. Steve Yegge noted that, as opposed to functional programming:
Rich Hickey, creator of Clojure, described object systems as overly simplistic models of the real world. He emphasized the inability of OOP to model time properly, which is getting increasingly problematic as software systems become more concurrent.
Alexander Stepanov compares object orientation unfavourably to generic programming:
Inheritance
OOP languages typically allow inheritance for code reuse and extensibility in the form of either classes or prototypes. These forms of inheritance are significantly different, but analogous terminology is used to define the concepts of object and instance.
Class-based
In class-based programming, the most popular style, each object is required to be an instance of a particular class. The class defines the data format or type (including member variables and their types) and available procedures (class methods or member functions) for a given type or class of object. Objects are created by calling a special type of method in the class known as a constructor. Classes may inherit from other classes, so they are arranged in a hierarchy that represents "is-a-type-of" relationships. For example, class Employee might inherit from class Person. All the data and methods available to the parent class also appear in the child class with the same names. For example, class Person might define variables "first_name" and "last_name" with method "make_full_name()". These will also be available in class Employee, which might add the variables "position" and "salary". It is guaranteed that all instances of class Employee will have the same variables, such as the name, position, and salary. Procedures and variables can be specific to either the class or the instance; this leads to the following terms:
Class variables – belong to the class as a whole; there is only one copy of each variable, shared across all instances of the class
Instance variables or attributes – data that belongs to individual objects; every object has its own copy of each one. All 4 variables mentioned above (first_name, position etc) are instance variables.
Member variables – refers to both the class and instance variables that are defined by a particular class.
Class methods – belong to the class as a whole and have access to only class variables and inputs from the procedure call
Instance methods – belong to individual objects, and have access to instance variables for the specific object they are called on, inputs, and class variables
Depending on the definition of the language, subclasses may or may not be able to override the methods defined by superclasses. Multiple inheritance is allowed in some languages, though this can make resolving overrides complicated. Some languages have special support for other concepts like traits and mixins, though, in any language with multiple inheritance, a mixin is simply a class that does not represent an is-a-type-of relationship. Mixins are typically used to add the same methods to multiple classes. For example, class UnicodeConversionMixin might provide a method unicode_to_ascii() when included in class FileReader and class WebPageScraper, which do not share a common parent.
Abstract classes cannot be instantiated into objects; they exist only for inheritance into other "concrete" classes that can be instantiated. In Java, the final keyword can be used to prevent a class from being subclassed.
Prototype-based
In contrast, in prototype-based programming, objects are the primary entities. Generally, the concept of a "class" does not even exist. Rather, the prototype or parent of an object is just another object to which the object is linked. In Self, an object may have multiple or no parents, but in the most popular prototype-based language, Javascript, every object has one prototype link (and only one). New objects can be created based on already existing objects chosen as their prototype. You may call two different objects apple and orange a fruit if the object fruit exists, and both apple and orange have fruit as their prototype. The idea of the fruit class does not exist explicitly, but can be modeled as the equivalence class of the objects sharing the same prototype, or as the set of objects satisfying a certain interface (duck typing). Unlike class-based programming, it is typically possible in prototype-based languages to define attributes and methods not shared with other objects; for example, the attribute sugar_content may be defined in apple but not orange.
Absence
Some languages like Go do not support inheritance at all. Go states that it is object-oriented, and Bjarne Stroustrup, author of C++, has stated that it is possible to do OOP without inheritance. The doctrine of composition over inheritance advocates implementing has-a relationships using composition instead of inheritance. For example, instead of inheriting from class Person, class Employee could give each Employee object an internal Person object, which it then has the opportunity to hide from external code even if class Person has many public attributes or methods. Delegation is another language feature that can be used as an alternative to inheritance.
Rob Pike has criticized the OO mindset for preferring a multilevel type hierarchy with layered abstractions to a three-line lookup table. He has called object-oriented programming "the Roman numerals of computing".
Bob Martin states that because they are software, related classes do not necessarily share the relationships of the things they represent.
Dynamic dispatch/message passing
It is the responsibility of the object, not any external code, to select the procedural code to execute in response to a method call, typically by looking up the method at run time in a table associated with the object. This feature is known as dynamic dispatch. If the call variability relies on more than the single type of the object on which it is called (i.e. at least one other parameter object is involved in the method choice), one speaks of multiple dispatch. A method call is also known as message passing. It is conceptualized as a message (the name of the method and its input parameters) being passed to the object for dispatch.
Dispatch interacts with inheritance; if a method is not present in a given object or class, the dispatch is delegated to its parent object or class, and so on, going up the chain of inheritance.
Data abstraction and encapsulation
Data abstraction is a design pattern in which data are visible only to semantically related functions, to prevent misuse. The success of data abstraction leads to frequent incorporation of data hiding as a design principle in object-oriented and pure functional programming. Similarly, encapsulation prevents external code from being concerned with the internal workings of an object. This facilitates code refactoring, for example allowing the author of the class to change how objects of that class represent their data internally without changing any external code (as long as "public" method calls work the same way). It also encourages programmers to put all the code that is concerned with a certain set of data in the same class, which organizes it for easy comprehension by other programmers. Encapsulation is a technique that encourages decoupling.
In object oriented programming, objects provide a layer which can be used to separate internal from external code and implement abstraction and encapsulation. External code can only use an object by calling a specific instance method with a certain set of input parameters, reading an instance variable, or writing to an instance variable. A program may create many instances of objects as it runs, which operate independently. This technique, it is claimed, allows easy re-use of the same procedures and data definitions for different sets of data, in addition to potentially mirroring real-world relationships intuitively. Rather than utilizing database tables and programming subroutines, the developer utilizes objects the user may be more familiar with: objects from their application domain. These claims that the OOP paradigm enhances reusability and modularity have been criticized.
The initial design is encouraged to use the most restrictive visibility possible, in order of local (or method) variables, private variables (in object oriented programming), and global (or public) variables, and only be expanded when and as much as necessary. This prevents changes to visibility from invalidating existing code.
If a class does not allow calling code to access internal object data and permits access through methods only, this is also a form of information hiding. Some languages (Java, for example) let classes enforce access restrictions explicitly, for example, denoting internal data with the private keyword and designating methods intended for use by code outside the class with the public keyword. Methods may also be designed public, private, or intermediate levels such as protected (which allows access from the same class and its subclasses, but not objects of a different class). In other languages (like Python) this is enforced only by convention (for example, private methods may have names that start with an underscore). In C#, Swift & Kotlin languages, internal keyword permits access only to files present in the same assembly, package, or module as that of the class.
In programming languages, particularly object-oriented ones, the emphasis on abstraction is vital. Object-oriented languages extend the notion of type to incorporate data abstraction, highlighting the significance of restricting access to internal data through methods. Eric S. Raymond has written that object-oriented programming languages tend to encourage thickly layered programs that destroy transparency. Raymond compares this unfavourably to the approach taken with Unix and the C programming language.
The "open/closed principle" advocates that classes and functions "should be open for extension, but closed for modification". Luca Cardelli has claimed that OOP languages have "extremely poor modularity properties with respect to class extension and modification", and tend to be extremely complex. The latter point is reiterated by Joe Armstrong, the principal inventor of Erlang, who is quoted as saying:
Leo Brodie has suggested a connection between the standalone nature of objects and a tendency to duplicate code in violation of the don't repeat yourself principle of software development.
Polymorphism
Subtyping – a form of polymorphism – is when calling code can be independent of which class in the supported hierarchy it is operating on – the parent class or one of its descendants. Meanwhile, the same operation name among objects in an inheritance hierarchy may behave differently.
For example, objects of the type Circle and Square are derived from a common class called Shape. The Draw function for each type of Shape implements what is necessary to draw itself while calling code can remain indifferent to the particular type of Shape being drawn.
This is another type of abstraction that simplifies code external to the class hierarchy and enables strong separation of concerns.
Open recursion
A common feature of objects is that methods are attached to them and can access and modify the object's data fields. In this brand of OOP, there is usually a special name such as or used to refer to the current object. In languages that support open recursion, object methods can call other methods on the same object (including themselves) using this name. This variable is late-bound; it allows a method defined in one class to invoke another method that is defined later, in some subclass thereof.
OOP languages
Simula (1967) is generally accepted as being the first language with the primary features of an object-oriented language. It was created for making simulation programs, in which what came to be called objects were the most important information representation. Smalltalk (1972 to 1980) is another early example and the one with which much of the theory of OOP was developed. Concerning the degree of object orientation, the following distinctions can be made:
Languages called "pure" OO languages, because everything in them is treated consistently as an object, from primitives such as characters and punctuation, all the way up to whole classes, prototypes, blocks, modules, etc. They were designed specifically to facilitate, even enforce, OO methods. Examples: Ruby, Scala, Smalltalk, Eiffel, Emerald, JADE, Self, Raku.
Languages designed mainly for OO programming, but with some procedural elements. Examples: Java, Python, C++, C#, Delphi/Object Pascal, VB.NET.
Languages that are historically procedural languages, but have been extended with some OO features. Examples: PHP, JavaScript, Perl, Visual Basic (derived from BASIC), MATLAB, COBOL 2002, Fortran 2003, ABAP, Ada 95, Pascal.
Languages with most of the features of objects (classes, methods, inheritance), but in a distinctly original form. Examples: Oberon (Oberon-1 or Oberon-2).
Languages with abstract data type support which may be used to resemble OO programming, but without all features of object-orientation. This includes object-based and prototype-based languages. Examples: JavaScript, Lua, Modula-2, CLU.
Chameleon languages that support multiple paradigms, including OO. Tcl stands out among these for TclOO, a hybrid object system that supports both prototype-based programming and class-based OO.
Popularity and reception
Many widely used languages, such as C++, Java, and Python, provide object-oriented features. Although in the past object-oriented programming was widely accepted, more recently essays criticizing object-oriented programming and recommending the avoidance of these features (generally in favor of functional programming) have been very popular in the developer community. Paul Graham has suggested that OOP's popularity within large companies is due to "large (and frequently changing) groups of mediocre programmers". According to Graham, the discipline imposed by OOP prevents any one programmer from "doing too much damage". Eric S. Raymond, a Unix programmer and open-source software advocate, has been critical of claims that present object-oriented programming as the "One True Solution".
Richard Feldman argues that these languages may have improved their modularity by adding OO features, but they became popular for reasons other than being object-oriented. In an article, Lawrence Krubner claimed that compared to other languages (LISP dialects, functional languages, etc.) OOP languages have no unique strengths, and inflict a heavy burden of unneeded complexity. A study by Potok et al. has shown no significant difference in productivity between OOP and procedural approaches. Luca Cardelli has claimed that OOP code is "intrinsically less efficient" than procedural code and that OOP can take longer to compile.
OOP in dynamic languages
In recent years, object-oriented programming has become especially popular in dynamic programming languages. Python, PowerShell, Ruby and Groovy are dynamic languages built on OOP principles, while Perl and PHP have been adding object-oriented features since Perl 5 and PHP 4, and ColdFusion since version 6.
The Document Object Model of HTML, XHTML, and XML documents on the Internet has bindings to the popular JavaScript/ECMAScript language. JavaScript is perhaps the best known prototype-based programming language, which employs cloning from prototypes rather than inheriting from a class (contrast to class-based programming). Another scripting language that takes this approach is Lua.
OOP in a network protocol
The messages that flow between computers to request services in a client-server environment can be designed as the linearizations of objects defined by class objects known to both the client and the server. For example, a simple linearized object would consist of a length field, a code point identifying the class, and a data value. A more complex example would be a command consisting of the length and code point of the command and values consisting of linearized objects representing the command's parameters. Each such command must be directed by the server to an object whose class (or superclass) recognizes the command and can provide the requested service. Clients and servers are best modeled as complex object-oriented structures. Distributed Data Management Architecture (DDM) took this approach and used class objects to define objects at four levels of a formal hierarchy:
Fields defining the data values that form messages, such as their length, code point and data values.
Objects and collections of objects similar to what would be found in a Smalltalk program for messages and parameters.
Managers similar to IBM i Objects, such as a directory to files and files consisting of metadata and records. Managers conceptually provide memory and processing resources for their contained objects.
A client or server consisting of all the managers necessary to implement a full processing environment, supporting such aspects as directory services, security, and concurrency control.
The initial version of DDM defined distributed file services. It was later extended to be the foundation of Distributed Relational Database Architecture (DRDA).
Design patterns
One way to address challenges of object-oriented design is via design patterns which are solution patterns to commonly occurring problems in software design. Some of these commonly occurring problems have implications and solutions particular to object-oriented development.
Object patterns
The following are notable software design patterns for OOP objects.
Function object: with a single method (in C++, the function operator, operator()) it acts much like a function
Immutable object: does not change state after creation
First-class object: can be used without restriction
Container object: contains other objects
Factory object: creates other objects
Metaobject: from which other objects can be created (compare with a class, which is not necessarily an object)
Prototype object: a specialized metaobject from which other objects can be created by copying
Singleton object: only instance of its class for the lifetime of the program
Filter object: receives a stream of data as its input and transforms it into the object's output
As an example of an object anti-pattern, the God object knows or does too much.
Inheritance and behavioral subtyping
It is intuitive to assume that inheritance creates a semantic "is a" relationship, and thus to infer that objects instantiated from subclasses can always be safely used instead of those instantiated from the superclass. This intuition is unfortunately false in most OOP languages, in particular in all those that allow mutable objects. Subtype polymorphism as enforced by the type checker in OOP languages (with mutable objects) cannot guarantee behavioral subtyping in any context. Behavioral subtyping is undecidable in general, so it cannot be implemented by a program (compiler). Class or object hierarchies must be carefully designed, considering possible incorrect uses that cannot be detected syntactically. This issue is known as the Liskov substitution principle.
Gang of Four design patterns
Design Patterns: Elements of Reusable Object-Oriented Software is an influential book published in 1994 by Erich Gamma, Richard Helm, Ralph Johnson, and John Vlissides, often referred to humorously as the "Gang of Four". Along with exploring the capabilities and pitfalls of object-oriented programming, it describes 23 common programming problems and patterns for solving them.
The book describes the following patterns:
Creational patterns (5): Factory method pattern, Abstract factory pattern, Singleton pattern, Builder pattern, Prototype pattern
Structural patterns (7): Adapter pattern, Bridge pattern, Composite pattern, Decorator pattern, Facade pattern, Flyweight pattern, Proxy pattern
Behavioral patterns (11): Chain-of-responsibility pattern, Command pattern, Interpreter pattern, Iterator pattern, Mediator pattern, Memento pattern, Observer pattern, State pattern, Strategy pattern, Template method pattern, Visitor pattern
Object-orientation and databases
Both object-oriented programming and relational database management systems (RDBMSs) are extremely common in software . Since relational databases do not store objects directly (though some RDBMSs have object-oriented features to approximate this), there is a general need to bridge the two worlds. The problem of bridging object-oriented programming accesses and data patterns with relational databases is known as object-relational impedance mismatch. There are some approaches to cope with this problem, but no general solution without downsides. One of the most common approaches is object-relational mapping, as found in IDE languages such as Visual FoxPro and libraries such as Java Data Objects and Ruby on Rails' ActiveRecord.
There are also object databases that can be used to replace RDBMSs, but these have not been as technically and commercially successful as RDBMSs.
Real-world modeling and relationships
OOP can be used to associate real-world objects and processes with digital counterparts. However, not everyone agrees that OOP facilitates direct real-world mapping or that real-world mapping is even a worthy goal; Bertrand Meyer argues in Object-Oriented Software Construction that a program is not a model of the world but a model of some part of the world; "Reality is a cousin twice removed". At the same time, some principal limitations of OOP have been noted.
For example, the circle-ellipse problem is difficult to handle using OOP's concept of inheritance.
However, Niklaus Wirth (who popularized the adage now known as Wirth's law: "Software is getting slower more rapidly than hardware becomes faster") said of OOP in his paper, "Good Ideas through the Looking Glass", "This paradigm closely reflects the structure of systems in the real world and is therefore well suited to model complex systems with complex behavior"
(contrast KISS principle).
Steve Yegge and others noted that natural languages lack the OOP approach of strictly prioritizing things (objects/nouns) before actions (methods/verbs). This problem may cause OOP to suffer more convoluted solutions than procedural programming.
OOP and control flow
OOP was developed to increase the reusability and maintainability of source code. Transparent representation of the control flow had no priority and was meant to be handled by a compiler. With the increasing relevance of parallel hardware and multithreaded coding, developing transparent control flow becomes more important, something hard to achieve with OOP.
Responsibility- vs. data-driven design
Responsibility-driven design defines classes in terms of a contract, that is, a class should be defined around a responsibility and the information that it shares. This is contrasted by Wirfs-Brock and Wilkerson with data-driven design, where classes are defined around the data-structures that must be held. The authors hold that responsibility-driven design is preferable.
SOLID and GRASP guidelines
SOLID is a mnemonic invented by Michael Feathers which spells out five software engineering design principles:
Single responsibility principle
Open/closed principle
Liskov substitution principle
Interface segregation principle
Dependency inversion principle
GRASP (General Responsibility Assignment Software Patterns) is another set of guidelines advocated by Craig Larman.
Formal semantics
Objects are the run-time entities in an object-oriented system. They may represent a person, a place, a bank account, a table of data, or any item that the program has to handle.
There have been several attempts at formalizing the concepts used in object-oriented programming. The following concepts and constructs have been used as interpretations of OOP concepts:
co algebraic data types
recursive types
encapsulated state
inheritance
records are the basis for understanding objects if function literals can be stored in fields (like in functional-programming languages), but the actual calculi need be considerably more complex to incorporate essential features of OOP. Several extensions of System F<: that deal with mutable objects have been studied; these allow both subtype polymorphism and parametric polymorphism (generics)
Attempts to find a consensus definition or theory behind objects have not proven very successful (however, see Abadi & Cardelli, A Theory of Objects for formal definitions of many OOP concepts and constructs), and often diverge widely. For example, some definitions focus on mental activities, and some on program structuring. One of the simpler definitions is that OOP is the act of using "map" data structures or arrays that can contain functions and pointers to other maps, all with some syntactic and scoping sugar on top. Inheritance can be performed by cloning the maps (sometimes called "prototyping").
Systems
CADES
Common Object Request Broker Architecture (CORBA)
Distributed Component Object Model
Distributed Data Management Architecture
Jeroo
Modeling languages
IDEF4
Interface description language
UML
| Technology | Programming | null |
3730679 | https://en.wikipedia.org/wiki/Chalicotherium | Chalicotherium | Chalicotherium (Ancient Greek /, -: pebble/gravel + /, diminutive of / : beast) is a genus of extinct odd-toed ungulates of the order Perissodactyla and family Chalicotheriidae. The genus is known from Europe and Asia, from the Middle Miocene to Late Miocene.
This animal would look much like other chalicotheriid species: an odd-looking herbivore with long clawed forelimbs and stouter weight-bearing hindlimbs.
The type species, Chalicotherium goldfussi, from Late Miocene Europe, was described by Johann Jakob Kaup in 1833. When the French naturalist George Cuvier first received a cleft claw from Eppelheim, Germany, he identified it as the toe bone of a gigantic pangolin.
Description
Chalicotherium, like many members of Perissodactyla, was adapted to browsing, though the chalicotheres were uniquely adapted to do so among ungulates. Its arms were long and heavily clawed, allowing them to walk on their knuckles only. The arms were used to reach for the branches of large trees and bring them close to its long head to strip them clean of leaves. The horse-like head itself shows adaptation to a diet of soft vegetation, since, as the animal reached sexual maturity, the incisors and upper canines were shed, suggesting that its muscular lips and the resulting gum pads were enough to crop fodder which was then processed by squarish, low-crowned molars.
Callosities on the ischium imply that these animals would sit on their haunches for extended periods of time, probably while feeding. Pad-supporting bony growth on the dorsal side of the manual phalanges is interpreted as evidence of knuckle-walking, which would probably be useful to avoid wearing down the claws, preserving them for use either as a forage-collecting rake or as a formidable defensive weapon, or both.
All of these characteristics show some convergence with such other creatures as ground sloths, great apes, bears (especially giant pandas), and a group of theropod dinosaurs known as therizinosaurs.
Classification
Taxonomic history
The type specimens for Chalicotherium goldfussi were found in the Upper Miocene strata of the Dinotherien-sande beds near Eppelsheim, in the Grand Duchy of Hesse, Germany. Johann Jakob Kaup, when describing this new animal in 1833, found the teeth to be pebble-like and named the creature accordingly. Later on, limbs found in strata located at Sansan in the department of Gers, Southwestern France, were first described as Macrotherium by Édouard Lartet in 1837. Further study of these fossil remains and subsequent finds by Filhol warranted a referral of the material described as Macrotherium to Chalicotherium.
Referral history for each species is detailed in the species list below along with morphological and geographical data where available.
Species
Valid:
Chalicotherium goldfussi J. J. Kaup, 1833.
The type species, it was found in Upper Miocene beds located in Germany. It weighed around 1500 kg and was 2.6 m high at the shoulder.
Chalicotherium brevirostris Colbert, 1934
First described as Macrotherium brevirostris, this species hails from the Upper Miocene Tung Gur Formation, Inner Mongolia, China.
Chalicotherium salinum Pickford, 1982
First described as Macrotherium salinum by Forster Cooper, this species was first discovered at the Lower Pliocene Lower Siwaliks beds in India; its chronological and geographic range was later extended to the Middle and Upper Miocene, and to Pakistan and China, respectively.
Invalid:
Chalicotherium antiquum J. J. Kaup, 1833.
Found at the same locality as the type species, it was later found wanting of diagnostic features and sunk into the type species.
Misassigned specimens:
Chalicotherium cf. C. brevirostris Wang et al., 2001.
Hailing for the Tsaidam Basin, northern Qinghai-Tibetan Plateau, China.
"Chalicotherium modicum" Stehlin, 1905.
A nomen nudum, actually a Schizotherium priscum tooth.
"Chalicotherium" bilobatum Cope.
Hailing from the Oligocene of Saskatchewan, this very fragmentary specimen was the type on which Russel erected the genus Oreinotherium.
Chalicotherium spp.
Specimens found in two Tajikistan localities, thought to pertain to at least two different species.
| Biology and health sciences | Perissodactyla | Animals |
3734085 | https://en.wikipedia.org/wiki/Animal%20breeding | Animal breeding | Animal breeding is a branch of animal science that addresses the evaluation (using best linear unbiased prediction and other methods) of the genetic value (estimated breeding value, EBV) of livestock. Selecting for breeding animals with superior EBV in growth rate, egg, meat, milk, or wool production, or with other desirable traits has revolutionized livestock production throughout the entire world. The scientific theory of animal breeding incorporates population genetics, quantitative genetics, statistics, and recently molecular genetics and is based on the pioneering work of Sewall Wright, Jay Lush, and Charles Henderson.
Breeding stock
Breeding stock is a group of animals used for the purpose of planned breeding. When individuals are looking to breed animals, they look for certain valuable traits in purebred animals, or may intend to use some type of crossbreeding to produce a new type of stock with different, and presumably superior abilities in a given area of endeavor. For example, when breeding swine for meat, the "breeding stock should be sound, fast growing, muscular, lean, and reproductively efficient." The "subjective selection of breeding stock" in horses has led to many horse breeds with particular performance traits. While breeding animals is common in an agricultural setting, it is also a common practice for the purpose of selling animals meant as pets, such as cats, dogs, horses, and birds, as well as less common animals, such as reptiles or some primates.
Purebred breeding
Mating animals of the same breed for maintaining such breed is referred to as purebred breeding. Opposite to the practice of mating animals of different breeds, purebred breeding aims to establish and maintain stable traits, that animals will pass to the next generation. By "breeding the best to the best", employing a certain degree of inbreeding, considerable culling, and selection for "superior" qualities, one could develop a bloodline or "breed" superior in certain respects to the original base stock.
Such animals can be recorded with a breed registry, the organisation that maintains pedigrees and/or stud books.
The observable phenomenon of hybrid vigor stands in contrast to the notion of breed purity.
For laboratory purposes, organisms such as mice have been inbred to 100% pure lines, as offered for sale by the Jackson laboratory. But this is highly unusual and difficult to do for most organisms, in whose populations all individuals harbor recessive, deleterious gene variants (alleles).
| Technology | Animal husbandry | null |
3735922 | https://en.wikipedia.org/wiki/Basa%20%28fish%29 | Basa (fish) | Basa (Pangasius bocourti), as it is commonly referred to, is a species of primarily freshwater-dwelling catfish in the shark-catfish family, Pangasiidae, native to the Mekong and Chao Phraya river basins of Mainland Southeast Asia. Economically, these fish are important as a regional food source, and are also prized on the international market. Outside of Asia, such as in North America or Australia, they are often referred to as "basa fish" or "swai" or by their specific name, "bocourti". In the United Kingdom, all species of Pangasius may, legally, be described as "river cobbler", "cobbler", "basa", "pangasius" or simply "panga", as well as any of these names with the addition of "catfish". In the rest of mainland Europe, these fish are mostly sold as "pangasius" or "panga". In Asian fish markets, names for basa also include "Pacific dory" and "patin". Other, related shark-catfish species may occasionally be labeled—albeit incorrectly—as basa, including the iridescent shark (P. hypophthalmus) and the yellowtail catfish (P. pangasius).
Description
The body of the basa is stout and heavy. The rounded head is broader than it is long, with the blunt snout having a white band on its muzzle. This species grows to a maximum length of .
Ecology
Basa fish feed on plants. They spawn at the onset of flood season and the young are first seen in June, averaging about by mid-June.
Market
Some bogey fish are labelled as swai; they are often mislabelled as tonguefish in China.
"Catfish war" in the U.S.
In 2002, the United States accused Vietnam of dumping catfish, namely P. bocourti and P. hypophthalmus, on the American market, arguing that the Vietnamese exporters, who are subsidised by Vietnam's government, were engaged in unfair competition.
With pressure from the U.S. catfish industry, the United States Congress passed a law in 2003 preventing the imported fish from being labelled as catfish, as well as imposing additional tariffs on the imported fish. Under the U.S. Food and Drug Administration ruling, only species from the family Ictaluridae can be sold as true catfish. As a result, the Vietnamese exporters of this fish now label their products sold in the U.S. as basa fish, striped pangasius, swai or bocourti.
At the height of the "catfish war", U.S. catfish farmers and others were describing the imported catfish as an inferior product. However, Mississippi State University researchers found imported basa were preferred three-to-one to US catfish in a small (58 testers) blind taste test.
United Kingdom
Basa has become common in the UK as "Vietnamese river cobbler", "river cobbler", or "basa". It is mainly sold by large supermarkets, in both fresh and frozen forms, as a cheaper alternative to popular white fish such as cod or haddock. Young's uses it in some of its frozen fish products, under the name basa. The import of basa is subject to the same stringent EU regulations as other food imports, as set out in the CBI pangasius product fact sheet UK Trading Standards officers said that cobbler was being fraudulently sold as cod by some fish-and-chip retailers to take advantage of the much lower price of cobbler, which was about half that of cod. This practice was highlighted by the successful prosecution of two retailers, using DNA evidence, in 2009 and 2010. Sometimes pangasius is described, legally, simply as "fish", as in "fish and chips".
Environmental and health concerns
Several environmental organisations specialising in marine ecosystems have raised concerns surrounding basa; OceanWise, an environmental group associated with Canada's Vancouver Aquarium, has flagged farmed basa for its potential pollution of ecosystems and interference with wild species. The group stated: "Open cage farming in Southeast Asia is associated with disease transfer to wild basa. There are also concerns about feed quality, farm operating standards and the biological impact of using wild stock for culturing." The Monterey Bay Aquarium in California currently lists the species in its "red flag" or "avoid" category. Both groups cite USA-farmed catfish as a more sustainable alternative.
Tests by Asda and Tesco supermarket corporations in the UK have found no trace of toxic contaminants. Testing by the Australian Quarantine and Inspection Service found trace levels of malachite green, but no other contaminants; this was likely the result of antiparasitic treatment administered to the fish, such as in the treatment of ich (white-spot disease), a common affliction of captive fishes with which malachite green (or methylene blue) is often remedied.
One case has been reported of a person without a general fish allergy having an anaphylactic reaction to pangasius.
| Biology and health sciences | Siluriformes | Animals |
25631460 | https://en.wikipedia.org/wiki/Tooth | Tooth | A tooth (: teeth) is a hard, calcified structure found in the jaws (or mouths) of many vertebrates and used to break down food. Some animals, particularly carnivores and omnivores, also use teeth to help with capturing or wounding prey, tearing food, for defensive purposes, to intimidate other animals often including their own, or to carry prey or their young. The roots of teeth are covered by gums. Teeth are not made of bone, but rather of multiple tissues of varying density and hardness that originate from the outermost embryonic germ layer, the ectoderm.
The general structure of teeth is similar across the vertebrates, although there is considerable variation in their form and position. The teeth of mammals have deep roots, and this pattern is also found in some fish, and in crocodilians. In most teleost fish, however, the teeth are attached to the outer surface of the bone, while in lizards they are attached to the inner surface of the jaw by one side. In cartilaginous fish, such as sharks, the teeth are attached by tough ligaments to the hoops of cartilage that form the jaw.
Monophyodonts are animals that develop only one set of teeth, while diphyodonts grow an early set of deciduous teeth and a later set of permanent or "adult" teeth. Polyphyodonts grow many sets of teeth. For example, sharks, grow a new set of teeth every two weeks to replace worn teeth. Most extant mammals including humans are diphyodonts, but there are exceptions including elephants, kangaroos, and manatees, all of which are polyphyodonts.
Rodent incisors grow and wear away continually through gnawing, which helps maintain relatively constant length. The industry of the beaver is due in part to this qualification. Some rodents, such as voles and guinea pigs (but not mice), as well as lagomorpha (rabbits, hares and pikas), have continuously growing molars in addition to incisors. Also, tusks (in tusked mammals) grow almost throughout life.
Teeth are not always attached to the jaw, as they are in mammals. In many reptiles and fish, teeth are attached to the palate or to the floor of the mouth, forming additional rows inside those on the jaws proper. Some teleosts even have teeth in the pharynx. While not true teeth in the usual sense, the dermal denticles of sharks are almost identical in structure and are likely to have the same evolutionary origin. Indeed, teeth appear to have first evolved in sharks, and are not found in the more primitive jawless fish – while lampreys do have tooth-like structures on the tongue, these are in fact, composed of keratin, not of dentine or enamel, and bear no relationship to true teeth. Though "modern" teeth-like structures with dentine and enamel have been found in late conodonts, they are now supposed to have evolved independently of later vertebrates' teeth.
Living amphibians typically have small teeth, or none at all, since they commonly feed only on soft foods. In reptiles, teeth are generally simple and conical in shape, although there is some variation between species, most notably the venom-injecting fangs of snakes. The pattern of incisors, canines, premolars and molars is found only in mammals, and to varying extents, in their evolutionary ancestors. The numbers of these types of teeth vary greatly between species; zoologists use a standardised dental formula to describe the precise pattern in any given group.
Etymology
The word tooth comes from Proto-Germanic , derived from the Proto-Indo-European , which was composed of the root plus the active participle suffix , therefore literally meaning .
The irregular plural form teeth is the result of Germanic umlaut whereby vowels immediately preceding a high vocalic in the following syllable were raised. As the nominative plural ending of the Proto-Germanic consonant stems (to which belonged) was , the root vowel in the plural form (changed by this point to via unrelated phonological processes) was raised to /œː/, and later unrounded to /eː/, resulting in the alternation attested from Old English. Cf. also Old English and , from Proto-Germanic and respectively.
Cognate with Latin , Greek (), and Sanskrit .
Origin
Teeth are assumed to have evolved either from ectoderm denticles (scales, much like those on the skin of sharks) that folded and integrated into the mouth (called the "outside–in" theory), or from endoderm pharyngeal teeth (primarily formed in the pharynx of jawless vertebrates) (the "inside–out" theory). In addition, there is another theory stating that neural crest gene regulatory network, and neural crest-derived ectomesenchyme are the key to generate teeth (with any epithelium, either ectoderm or endoderm).
The genes governing tooth development in mammals are homologous to those involved in the development of fish scales. Study of a tooth plate of a fossil of the extinct fish Romundina stellina showed that the teeth and scales were made of the same tissues, also found in mammal teeth, lending support to the theory that teeth evolved as a modification of scales.
Mammals
Teeth are among the most distinctive (and long-lasting) features of mammal species. Paleontologists use teeth to identify fossil species and determine their relationships. The shape of the animal's teeth are related to its diet. For example, plant matter is hard to digest, so herbivores have many molars for chewing and grinding. Carnivores, on the other hand, have canine teeth to kill prey and to tear meat.
Mammals, in general, are diphyodont, meaning that they develop two sets of teeth. In humans, the first set (the "baby", "milk", "primary" or "deciduous" set) normally starts to appear at about six months of age, although some babies are born with one or more visible teeth, known as neonatal teeth. Normal tooth eruption at about six months is known as teething and can be painful. Kangaroos, elephants, and manatees are unusual among mammals because they are polyphyodonts.
Aardvark
In aardvarks, teeth lack enamel and have many pulp tubules, hence the name of the order Tubulidentata.
Canines
In dogs, the teeth are less likely than humans to form dental cavities because of the very high pH of dog saliva, which prevents enamel from demineralizing. Sometimes called cuspids, these teeth are shaped like points (cusps) and are used for tearing and grasping food.
Cetaceans
Like human teeth, whale teeth have polyp-like protrusions located on the root surface of the tooth. These polyps are made of cementum in both species, but in human teeth, the protrusions are located on the outside of the root, while in whales the nodule is located on the inside of the pulp chamber. While the roots of human teeth are made of cementum on the outer surface, whales have cementum on the entire surface of the tooth with a very small layer of enamel at the tip. This small enamel layer is only seen in older whales where the cementum has been worn away to show the underlying enamel.
The toothed whale is a parvorder of the cetaceans characterized by having teeth. The teeth differ considerably among the species. They may be numerous, with some dolphins bearing over 100 teeth in their jaws. On the other hand, the narwhals have a giant unicorn-like tusk, which is a tooth containing millions of sensory pathways and used for sensing during feeding, navigation, and mating. It is the most neurologically complex tooth known. Beaked whales are almost toothless, with only bizarre teeth found in males. These teeth may be used for feeding but also for demonstrating aggression and showmanship.
Primates
In humans (and most other primates), there are usually 20 primary (also "baby" or "milk") teeth, and later up to 32 permanent teeth. Four of these 32 may be third molars or wisdom teeth, although these are not present in all adults, and may be removed surgically later in life.
Among primary teeth, 10 of them are usually found in the maxilla (i.e. upper jaw) and the other 10 in the mandible (i.e. lower jaw). Among permanent teeth, 16 are found in the maxilla and the other 16 in the mandible. Most of the teeth have uniquely distinguishing features.
Horse
An adult horse has between 36 and 44 teeth. The enamel and dentin layers of horse teeth are intertwined. All horses have 12 premolars, 12 molars, and 12 incisors. Generally, all male equines also have four canine teeth (called tushes) between the molars and incisors. However, few female horses (less than 28%) have canines, and those that do usually have only one or two, which many times are only partially erupted. A few horses have one to four wolf teeth, which are vestigial premolars, with most of those having only one or two. They are equally common in male and female horses and much more likely to be on the upper jaw. If present these can cause problems as they can interfere with the horse's bit contact. Therefore, wolf teeth are commonly removed.
Horse teeth can be used to estimate the animal's age. Between birth and five years, age can be closely estimated by observing the eruption pattern on milk teeth and then permanent teeth. By age five, all permanent teeth have usually erupted. The horse is then said to have a "full" mouth. After the age of five, age can only be conjectured by studying the wear patterns on the incisors, shape, the angle at which the incisors meet, and other factors. The wear of teeth may also be affected by diet, natural abnormalities, and cribbing. Two horses of the same age may have different wear patterns.
A horse's incisors, premolars, and molars, once fully developed, continue to erupt as the grinding surface is worn down through chewing. A young adult horse will have teeth, which are long, with the majority of the crown remaining below the gumline in the dental socket. The rest of the tooth will slowly emerge from the jaw, erupting about each year, as the horse ages. When the animal reaches old age, the crowns of the teeth are very short and the teeth are often lost altogether. Very old horses, if lacking molars, may need to have their fodder ground up and soaked in water to create a soft mush for them to eat in order to obtain adequate nutrition.
Proboscideans
Elephants' tusks are specialized incisors for digging food up and fighting. Some elephant teeth are similar to those in manatees, and elephants are believed to have undergone an aquatic phase in their evolution.
At birth, elephants have a total of 28 molar plate-like grinding teeth not including the tusks. These are organized into four sets of seven successively larger teeth which the elephant will slowly wear through during its lifetime of chewing rough plant material. Only four teeth are used for chewing at a given time, and as each tooth wears out, another tooth moves forward to take its place in a process similar to a conveyor belt. The last and largest of these teeth usually becomes exposed when the animal is around 40 years of age, and will often last for an additional 20 years. When the last of these teeth has fallen out, regardless of the elephant's age, the animal will no longer be able to chew food and will die of starvation.
Rabbit
Rabbits and other lagomorphs usually shed their deciduous teeth before (or very shortly after) their birth, and are usually born with their permanent teeth. The teeth of rabbits complement their diet, which consists of a wide range of vegetation. Since many of the foods are abrasive enough to cause attrition, rabbit teeth grow continuously throughout life. Rabbits have a total of six incisors, three upper premolars, three upper molars, two lower premolars, and two lower molars on each side. There are no canines. Dental formula is = 28. Three to four millimeters of the tooth is worn away by incisors every week, whereas the cheek teeth require a month to wear away the same amount.
The incisors and cheek teeth of rabbits are called aradicular hypsodont teeth. This is sometimes referred to as an elodent dentition. These teeth grow or erupt continuously. The growth or eruption is held in balance by dental abrasion from chewing a diet high in fiber.
Rodents
Rodents have upper and lower hypselodont incisors that can continuously grow enamel throughout its life without having properly formed roots. These teeth are also known as aradicular teeth, and unlike humans whose ameloblasts die after tooth development, rodents continually produce enamel, they must wear down their teeth by gnawing on various materials. Enamel and dentin are produced by the enamel organ, and growth is dependent on the presence of stem cells, cellular amplification, and cellular maturation structures in the odontogenic region. Rodent incisors are used for cutting wood, biting through the skin of fruit, or for defense. This allows for the rate of wear and tooth growth to be at equilibrium. The microstructure of rodent incisor enamel has shown to be useful in studying the phylogeny and systematics of rodents because of its independent evolution from the other dental traits. The enamel on rodent incisors are composed of two layers: the inner portio interna (PI) with Hunter-Schreger bands (HSB) and an outer portio externa (PE) with radial enamel (RE). It usually involves the differential regulation of the epithelial stem cell niche in the tooth of two rodent species, such as guinea pigs.
The teeth have enamel on the outside and exposed dentin on the inside, so they self-sharpen during gnawing. On the other hand, continually growing molars are found in some rodent species, such as the sibling vole and the guinea pig. There is variation in the dentition of the rodents, but generally, rodents lack canines and premolars, and have a space between their incisors and molars, called the diastema region.
Manatee
Manatees are polyphyodont with mandibular molars developing separately from the jaw and are encased in a bony shell separated by soft tissue.
Walrus
Walrus tusks are canine teeth that grow continuously throughout life.
Fish
Fish, such as sharks, may go through many teeth in their lifetime. The replacement of multiple teeth is known as polyphyodontia.
A class of prehistoric shark are called cladodonts for their strange forked teeth.
Unlike the continuous shedding of functional teeth seen in modern sharks, the majority of stem chondrichthyan lineages retained all tooth generations developed throughout the life of the animal. This replacement mechanism is exemplified by the tooth whorl-based dentitions of acanthodians, which include the oldest known toothed vertebrate, Qianodus duplicis.
Amphibians
All amphibians have pedicellate teeth, which are modified to be flexible due to connective tissue and uncalcified dentine that separates the crown from the base of the tooth.
Most amphibians exhibit teeth that have a slight attachment to the jaw or acrodont teeth. Acrodont teeth exhibit limited connection to the dentary and have little enervation. This is ideal for organisms who mostly use their teeth for grasping, but not for crushing and allows for rapid regeneration of teeth at a low energy cost. Teeth are usually lost in the course of feeding if the prey is struggling. Additionally, amphibians that undergo a metamorphosis develop bicuspid shaped teeth.
Reptiles
The teeth of reptiles are replaced constantly throughout their lives. Crocodilian juveniles replace teeth with larger ones at a rate as high as one new tooth per socket every month. Once mature, tooth replacement rates can slow to two years and even longer. Overall, crocodilians may use 3,000 teeth from birth to death. New teeth are created within old teeth.
Birds
A skull of Ichthyornis discovered in 2014 suggests that the beak of birds may have evolved from teeth to allow chicks to escape their shells earlier, and thus avoid predators and also to penetrate protective covers such as hard earth to access underlying food.
Invertebrates
True teeth are unique to vertebrates, although many invertebrates have analogous structures often referred to as teeth. The organisms with the simplest genome bearing such tooth-like structures are perhaps the parasitic worms of the family Ancylostomatidae. For example, the hookworm Necator americanus has two dorsal and two ventral cutting plates or teeth around the anterior margin of the buccal capsule. It also has a pair of subdorsal and a pair of subventral teeth located close to the rear.
Historically, the European medicinal leech, another invertebrate parasite, has been used in medicine to remove blood from patients. They have three jaws (tripartite) that resemble saws in both appearance and function, and on them are about 100 sharp teeth used to incise the host. The incision leaves a mark that is an inverted Y inside of a circle. After piercing the skin and injecting anticoagulants (hirudin) and anaesthetics, they suck out blood, consuming up to ten times their body weight in a single meal.
In some species of Bryozoa, the first part of the stomach forms a muscular gizzard lined with chitinous teeth that crush armoured prey such as diatoms. Wave-like peristaltic contractions then move the food through the stomach for digestion.
Molluscs have a structure called a radula, which bears a ribbon of chitinous teeth. However, these teeth are histologically and developmentally different from vertebrate teeth and are unlikely to be homologous. For example, vertebrate teeth develop from a neural crest mesenchyme-derived dental papilla, and the neural crest is specific to vertebrates, as are tissues such as enamel.
The radula is used by molluscs for feeding and is sometimes compared rather inaccurately to a tongue. It is a minutely toothed, chitinous ribbon, typically used for scraping or cutting food before the food enters the oesophagus. The radula is unique to molluscs, and is found in every class of mollusc apart from bivalves.
Within the gastropods, the radula is used in feeding by both herbivorous and carnivorous snails and slugs. The arrangement of teeth (also known as denticles) on the radula ribbon varies considerably from one group to another as shown in the diagram on the left.
Predatory marine snails such as the Naticidae use the radula plus an acidic secretion to bore through the shell of other molluscs. Other predatory marine snails, such as the Conidae, use a specialized radula tooth as a poisoned harpoon. Predatory pulmonate land slugs, such as the ghost slug, use elongated razor-sharp teeth on the radula to seize and devour earthworms. Predatory cephalopods, such as squid, use the radula for cutting prey.
In most of the more ancient lineages of gastropods, the radula is used to graze by scraping diatoms and other microscopic algae off rock surfaces and other substrates. Limpets scrape algae from rocks using radula equipped with exceptionally hard rasping teeth. These teeth have the strongest known tensile strength of any biological material, outperforming spider silk. The mineral protein of the limpet teeth can withstand a tensile stress of 4.9 GPa, compared to 4 GPa of spider silk and 0.5 GPa of human teeth.
Fossilization and taphonomy
Because teeth are very resistant, often preserved when bones are not, and reflect the diet of the host organism, they are very valuable to archaeologists and palaeontologists. Early fish such as the thelodonts had scales composed of dentine and an enamel-like compound, suggesting that the origin of teeth was from scales which were retained in the mouth. Fish as early as the late Cambrian had dentine in their exoskeletons, which may have functioned in defense or for sensing their environments. Dentine can be as hard as the rest of teeth and is composed of collagen fibres, reinforced with hydroxyapatite.
Though teeth are very resistant, they also can be brittle and highly susceptible to cracking. However, cracking of the tooth can be used as a diagnostic tool for predicting bite force. Additionally, enamel fractures can also give valuable insight into the diet and behaviour of archaeological and fossil samples.
Decalcification removes the enamel from teeth and leaves only the organic interior intact, which comprises dentine and cementine. Enamel is quickly decalcified in acids, perhaps by dissolution by plant acids or via diagenetic solutions, or in the stomachs of vertebrate predators. Enamel can be lost by abrasion or spalling, and is lost before dentine or bone are destroyed by the fossilisation process. In such a case, the 'skeleton' of the teeth would consist of the dentine, with a hollow pulp cavity.
The organic part of dentine, conversely, is destroyed by alkalis.
| Biology and health sciences | Gastrointestinal tract | Biology |
6668150 | https://en.wikipedia.org/wiki/Clearing%20the%20neighbourhood | Clearing the neighbourhood | In celestial mechanics, "clearing the neighbourhood" (or dynamical dominance) around a celestial body's orbit describes the body becoming gravitationally dominant such that there are no other bodies of comparable size other than its natural satellites or those otherwise under its gravitational influence.
"Clearing the neighbourhood" is one of three necessary criteria for a celestial body to be considered a planet in the Solar System, according to the definition adopted in 2006 by the International Astronomical Union (IAU). In 2015, a proposal was made to extend the definition to exoplanets.
In the end stages of planet formation, a planet, as so defined, will have "cleared the neighbourhood" of its own orbital zone, i.e. removed other bodies of comparable size. A large body that meets the other criteria for a planet but has not cleared its neighbourhood is classified as a dwarf planet. This includes Pluto, whose orbit is partly inside Neptune's and shares its orbital neighbourhood with many Kuiper belt objects. The IAU's definition does not attach specific numbers or equations to this term, but all IAU-recognised planets have cleared their neighbourhoods to a much greater extent (by orders of magnitude) than any dwarf planet or candidate for dwarf planet.
The phrase stems from a paper presented to the 2000 IAU general assembly by the planetary scientists Alan Stern and Harold F. Levison. The authors used several similar phrases as they developed a theoretical basis for determining if an object orbiting a star is likely to "clear its neighboring region" of planetesimals based on the object's mass and its orbital period. Steven Soter prefers to use the term dynamical dominance, and Jean-Luc Margot notes that such language "seems less prone to misinterpretation".
Prior to 2006, the IAU had no specific rules for naming planets, as no new planets had been discovered for decades, whereas there were well-established rules for naming an abundance of newly discovered small bodies such as asteroids or comets. The naming process for Eris stalled after the announcement of its discovery in 2005, because its size was comparable to that of Pluto. The IAU sought to resolve the naming of Eris by seeking a taxonomical definition to distinguish planets from minor planets.
Criteria
The phrase refers to an orbiting body (a planet or protoplanet) "sweeping out" its orbital region over time, by gravitationally interacting with smaller bodies nearby. Over many orbital cycles, a large body will tend to cause small bodies either to accrete with it, or to be disturbed to another orbit, or to be captured either as a satellite or into a resonant orbit. As a consequence it does not then share its orbital region with other bodies of significant size, except for its own satellites, or other bodies governed by its own gravitational influence. This latter restriction excludes objects whose orbits may cross but that will never collide with each other due to orbital resonance, such as Jupiter and its trojans, Earth and 3753 Cruithne, or Neptune and the plutinos. As to the extent of orbit clearing required, Jean-Luc Margot emphasises "a planet can never completely clear its orbital zone, because gravitational and radiative forces continually perturb the orbits of asteroids and comets into planet-crossing orbits" and states that the IAU did not intend the impossible standard of impeccable orbit clearing.
Stern–Levison's
In their paper, Stern and Levison sought an algorithm to determine which "planetary bodies control the region surrounding them". They defined (lambda), a measure of a body's ability to scatter smaller masses out of its orbital region over a period of time equal to the age of the Universe (Hubble time). is a dimensionless number defined as
where is the mass of the body, is the body's semi-major axis, and is a function of the orbital elements of the small body being scattered and the degree to which it must be scattered. In the domain of the solar planetary disc, there is little variation in the average values of for small bodies at a particular distance from the Sun.
If > 1, then the body will likely clear out the small bodies in its orbital zone. Stern and Levison used this discriminant to separate the gravitationally rounded, Sun-orbiting bodies into überplanets, which are "dynamically important enough to have cleared [their] neighboring planetesimals", and unterplanets. The überplanets are the eight most massive solar orbiters (i.e. the IAU planets), and the unterplanets are the rest (i.e. the IAU dwarf planets).
Soter's
Steven Soter proposed an observationally based measure (mu), which he called the "planetary discriminant", to separate bodies orbiting stars into planets and non-planets. He defines as
where is a dimensionless parameter, is the mass of the candidate planet, and is the mass of all other bodies that share an orbital zone, that is all bodies whose orbits cross a common radial distance from the primary, and whose non-resonant periods differ by less than an order of magnitude.
The order-of-magnitude similarity in period requirement excludes comets from the calculation, but the combined mass of the comets turns out to be negligible compared with the other small Solar System bodies, so their inclusion would have little impact on the results. μ is then calculated by dividing the mass of the candidate body by the total mass of the other objects that share its orbital zone. It is a measure of the actual degree of cleanliness of the orbital zone. Soter proposed that if > 100, then the candidate body be regarded as a planet.
Margot's
Astronomer Jean-Luc Margot has proposed a discriminant, (pi), that can categorise a body based only on its own mass, its semi-major axis, and its star's mass. Like Stern–Levison's , is a measure of the ability of the body to clear its orbit, but unlike , it is solely based on theory and does not use empirical data from the Solar System. is based on properties that are feasibly determinable even for exoplanetary bodies, unlike Soter's , which requires an accurate census of the orbital zone.
where is the mass of the candidate body in Earth masses, is its semi-major axis in AU, is the mass of the parent star in solar masses, and is a constant chosen so that > 1 for a body that can clear its orbital zone. depends on the extent of clearing desired and the time required to do so. Margot selected an extent of times the Hill radius and a time limit of the parent star's lifetime on the main sequence (which is a function of the mass of the star). Then, in the mentioned units and a main-sequence lifetime of 10 billion years, = 807. The body is a planet if > 1. The minimum mass necessary to clear the given orbit is given when = 1.
is based on a calculation of the number of orbits required for the candidate body to impart enough energy to a small body in a nearby orbit such that the smaller body is cleared out of the desired orbital extent. This is unlike , which uses an average of the clearing times required for a sample of asteroids in the asteroid belt, and is thus biased to that region of the Solar System. 's use of the main-sequence lifetime means that the body will eventually clear an orbit around the star; 's use of a Hubble time means that the star might disrupt its planetary system (e.g. by going nova) before the object is actually able to clear its orbit.
The formula for assumes a circular orbit. Its adaptation to elliptical orbits is left for future work, but Margot expects it to be the same as that of a circular orbit to within an order of magnitude.
To accommodate planets in orbit around brown dwarfs, an updated version of the criterion with a uniform clearing time scale of 10 billion years was published in 2024. The values of for Solar System bodies remain unchanged.
Numerical values
Below is a list of planets and dwarf planets ranked by Margot's planetary discriminant , in decreasing order. For all eight planets defined by the IAU, is orders of magnitude greater than 1, whereas for all dwarf planets, is orders of magnitude less than 1. Also listed are Stern–Levison's and Soter's ; again, the planets are orders of magnitude greater than 1 for and 100 for , and the dwarf planets are orders of magnitude less than 1 for and 100 for . Also shown are the distances where = 1 and = 1 (where the body would change from being a planet to being a dwarf planet).
The mass of Sedna is not known; it is very roughly estimated here as , on the assumption of a density of about .
Disagreement
Stern, the principal investigator of the New Horizons mission to Pluto, disagreed with the reclassification of Pluto on the basis of its inability to clear a neighbourhood. He argued that the IAU's wording is vague, and that — like Pluto — Earth, Mars, Jupiter and Neptune have not cleared their orbital neighbourhoods either. Earth co-orbits with 10,000 near-Earth asteroids (NEAs), and Jupiter has 100,000 trojans in its orbital path. "If Neptune had cleared its zone, Pluto wouldn't be there", he said.
The IAU category of 'planets' is nearly identical to Stern's own proposed category of 'überplanets'. In the paper proposing Stern and Levison's discriminant, they stated, "we define an überplanet as a planetary body in orbit about a star that is dynamically important enough to have cleared its neighboring planetesimals ..." and a few paragraphs later, "From a dynamical standpoint, our solar system clearly contains 8 überplanets" — including Earth, Mars, Jupiter, and Neptune. Although Stern proposed this to define dynamical subcategories of planets, he rejected it for defining what a planet is, advocating the use of intrinsic attributes over dynamical relationships.
| Physical sciences | Planetary science | Astronomy |
2774053 | https://en.wikipedia.org/wiki/Eudicots | Eudicots | The eudicots, Eudicotidae, or eudicotyledons are a clade of flowering plants (angiosperms) which are mainly characterized by having two seed leaves (cotyledons) upon germination. The term derives from dicotyledon (etymologically, eu = true; di = two; cotyledon = seed leaf). Historically, authors have used the terms tricolpates or non-magnoliid dicots. The current botanical terms were introduced in 1991, by evolutionary botanist James A. Doyle and paleobotanist Carol L. Hotton, to emphasize the later evolutionary divergence of tricolpate dicots from earlier, less specialized, dicots.
Scores of familiar plants are eudicots, including many commonly cultivated and edible plants, numerous trees, tropicals and ornamentals. Among the most well-known eudicot genera are those of the sunflower (Helianthus), dandelion (Taraxacum), forget-me-not (Myosotis), cabbage (Brassica), apple (Malus), buttercup (Ranunculus), maple (Acer) and macadamia (Macadamia). Most leafy, mid-latitude trees are also classified as eudicots, with notable exceptions being the magnolias and American tulip tree (Liriodendron)—which belong to the magnoliids—and Ginkgo biloba, which is not an angiosperm.
Description
The close relationships among flowering plants with tricolpate pollen grains was initially seen in morphological studies of shared derived characters. These plants have a distinct trait in their pollen grains of exhibiting three colpi or grooves paralleling the polar axis.
Later molecular evidence confirmed the genetic basis for the evolutionary relationships among flowering plants with tricolpate pollen grains and dicotyledonous traits. The term means "true dicotyledons", as it contains the majority of plants that have been considered dicots and have characteristics of the dicots. One of the genetic traits which defines the eudicots is the duplication of DELLA protein-encoding genes in their most recent common ancestor. The term "eudicots" has subsequently been widely adopted in botany to refer to one of the two largest clades of angiosperms (constituting over 70% of the angiosperm species), monocots being the other. The remaining angiosperms include magnoliids and what are sometimes referred to as basal angiosperms or paleodicots, but these terms have not been widely or consistently adopted, as they do not refer to a monophyletic group.
According to molecular clock calculations, the lineage that led to eudicots split from other plants about 134 million years ago or 155-160 million years ago.
Taxonomy
The earlier name for the eudicots is tricolpates, a name which refers to the grooved structure of the pollen. Members of the group have tricolpate pollen, or forms derived from it. These pollens have three or more pores set in furrows called colpi. In contrast, most of the other seed plants (that is the gymnosperms, the monocots and the paleodicots) produce monosulcate pollen, with a single pore set in a differently oriented groove called the sulcus. The name "tricolpates" is preferred by some botanists to avoid confusion with the dicots, a nonmonophyletic group.
The name "eudicots" (plural) is used in the APG systems (from APG system, of 1998, to APG IV system, of 2016) for classification of angiosperms. It is applied to a clade, a monophyletic group, which includes most of the (former) dicots.
"Tricolpate" is a synonym for the "Eudicot" monophyletic group, the "true dicotyledons" (which are distinguished from all other flowering plants by their tricolpate pollen structure). The number of pollen grain furrows or pores helps classify the flowering plants, with eudicots having three colpi (tricolpate), and other groups having one sulcus.
Pollen apertures are any modification of the wall of the pollen grain. These modifications include thinning, ridges and pores, they serve as an exit for the pollen contents and allow shrinking and swelling of the grain caused by changes in moisture content. The elongated apertures/ furrows in the pollen grain are called colpi (singular colpus), which, along with pores, are a chief criterion for identifying the pollen classes.
Subdivisions
The eudicots can be divided into two groups: the basal eudicots and the core eudicots. Basal eudicot is an informal name for a paraphyletic group. The core eudicots are a monophyletic group. A 2010 study suggested the core eudicots can be divided into two clades, Gunnerales and a clade called Pentapetalae, comprising all the remaining core eudicots.
The Pentapetalae can be then divided into three clades:
Dilleniales
superrosids consisting of Saxifragales and rosids (the APG IV system includes the Vitales in the rosids)
superasterids consisting of Santalales, Berberidopsidales, Caryophyllales and asterids
This division of the eudicots is shown in the following cladogram:
The following is a more detailed breakdown according to APG IV, showing within each clade and orders:
clade Eudicots
order Ranunculales
order Proteales
order Trochodendrales
order Buxales
clade Core eudicots
order Gunnerales
order Dilleniales
clade Superrosids
order Saxifragales
clade Rosids
order Vitales
clade Fabids
order Fabales
order Rosales
order Fagales
order Cucurbitales
order Oxalidales
order Malpighiales
order Celastrales
order Zygophyllales
clade Malvids
order Geraniales
order Myrtales
order Crossosomatales
order Picramniales
order Malvales
order Brassicales
order Huerteales
order Sapindales
clade Superasterids
order Berberidopsidales
order Santalales
order Caryophyllales
clade Asterids
order Cornales
order Ericales
clade Campanulids
order Aquifoliales
order Asterales
order Escalloniales
order Bruniales
order Apiales
order Dipsacales
order Paracryphiales
clade Lamiids
order Solanales
order Lamiales
order Vahliales
order Gentianales
order Boraginales
order Garryales
order Metteniusales
order Icacinales
| Biology and health sciences | Eudicots | null |
2776501 | https://en.wikipedia.org/wiki/Visible-light%20astronomy | Visible-light astronomy | Visible-light astronomy encompasses a wide variety of astronomical observation via telescopes that are sensitive in the range of visible light (optical telescopes). Visible-light astronomy is part of optical astronomy, and differs from astronomies based on invisible types of light in the electromagnetic radiation spectrum, such as radio waves, infrared waves, ultraviolet waves, X-ray waves and gamma-ray waves. Visible light ranges from 380 to 750 nanometers in wavelength.
Visible-light astronomy has existed as long as people have been looking up at the night sky, although it has since improved in its observational capabilities since the invention of the telescope, which is commonly credited to Hans Lippershey, a German-Dutch spectacle-maker, although Galileo played a large role in the development and creation of telescopes.
Since visible-light astronomy is restricted to only visible light, no equipment is necessary for simply star gazing. This means that it's the most commonly participated in type of astronomy, as well as the oldest.
History
Beginning
Before the advent of telescopes, astronomy was limited solely to unaided eyesight. Humans have been gazing at stars and other objects in the night sky for thousands of years, as is evident in the naming of many constellations, notably the largely Greek names used today.
Hans Lippershey, a German-Dutch spectacle maker, is commonly credited as being the first to invent the optical telescope. Lippershey is the first recorded person to apply for a patent for a telescope; however, it is unclear if Lippershey was the first to build a telescope. Based only on uncertain descriptions of the telescope for which Lippershey tried to obtain a patent, Galileo Galilei made a telescope with about 3× magnification in the following year. Galileo later made improved versions with up to 30× magnification. With a Galilean telescope, the observer could see magnified, upright images on Earth; it was what is commonly known as a terrestrial telescope or a spyglass. Galileo could also use it to observe the sky, and for a time was one of those who could construct telescopes good enough for that purpose. On 25 August 1609, Galileo demonstrated one of his early telescopes, with a magnification of up to 8 or 9, to Venetian lawmakers. Galileo's telescopes were also a profitable sideline, selling them to merchants who found them useful both at sea and as items of trade. He published his initial telescopic astronomical observations in March 1610 in a brief treatise titled Sidereus Nuncius (Starry Messenger).
The human eye, now with optical aid, remained the only image sensor until the advent of astrophotography in the 19th century.
Modern day
In the modern day, visible-light astronomy is still practiced by many amateur astronomers, especially since telescopes are much more widely available for the public, as compared to when they were first being invented. Government agencies, such as NASA, are very involved in the modern day research and observation of visible objects and celestial bodies. In the modern day, the highest quality pictures and data are obtained via space telescopes; telescopes that are outside of the Earth's atmosphere. This allows for much clearer observations, as the atmosphere is not disrupting the image and viewing quality of the telescope, meaning objects can be observed in much greater detail, and much more distant or low-light objects may be observed. Additionally, this means that observations are able to be made at any time, rather than only during the night.
Hubble Space Telescope
The Hubble Space Telescope is a space telescope created by NASA, and was launched into low Earth orbit in 1990. It is still in operation today. The Hubble Space Telescope's four main instruments observe in the near ultraviolet, visible, and near infrared spectra. Hubble's images are some of the most detailed images ever taken, leading to many breakthroughs in astrophysics, such as accurately determining the rate of expansion of the universe.
Optical telescopes
There are three main types of telescopes used in visible-light astronomy:
Refracting telescopes, which use lenses to form the image. Commonly used by amateur astronomers, especially for viewing brighter objects such as the Moon, and planets, due to lower cost and ease of usage.
Reflecting telescopes, which use mirrors to form the image. Commonly used for scientific purposes.
Catadioptric telescopes, which use a combination of lenses and mirrors to form the image; essentially a combination of refracting and reflecting telescopes.
Each type of telescope suffers from different types of aberration; refracting telescopes have chromatic aberration, which causes colors to be shown on edges separating light and dark parts of the image, where there should not be such colors. This is due to the lens being unable to focus all colors to the same convergence point. Reflecting telescopes suffer from several types of optical inaccuracies, such as off-axis aberrations near the edges of the field of view. Catadioptric telescopes vary in the types of optical inaccuracies present, as there are numerous catadioptric telescope designs.
Effect of ambient brightness
The visibility of celestial objects in the night sky is affected by light pollution, with the presence of the Moon in the night sky historically hindering astronomical observation by increasing the amount of ambient lighting. With the advent of artificial light sources, however, light pollution has been a growing problem for viewing the night sky. Special filters and modifications to light fixtures can help to alleviate this problem, but for the best views, both professional and amateur optical astronomers seek viewing sites located far from major urban areas. In order to avoid light pollution of Earth's sky, among other reasons, many telescopes are put outside of the Earth's atmosphere, where not only light pollution, but also atmospheric distortion and obscuration are minimized.
Commonly observed objects
The most commonly observed objects tend to be ones that do not require a telescope to view, such as the Moon, meteors, planets, constellations, and stars.
The Moon is a very commonly observed astronomical object, especially by amateur astronomers and skygazers. This is due to several reasons: the Moon is the brightest object in the night sky, the Moon is the largest object in the night sky, and the Moon has long been significant in many cultures, such as being the basis for many calendars. The Moon also does not require any kind of telescope or binoculars to see effectively, making it extremely convenient and common for people to observe.
Meteors, often called "shooting stars" are also commonly observed. Meteor showers, such as the Perseids and Leonids, make viewing meteors much easier, as a multitude of meteors are visible in a relatively short period of time.
Planets are usually observed with the aid of a telescope or binoculars. Venus is likely the easiest planet to observe without the aid of any instruments, as it is very bright, and can even be seen in daylight. However, Mars, Jupiter, and Saturn can also be seen without the aid of telescopes or binoculars.
Constellations and stars are also often observed, and have been used in the past for navigation, especially by ships at sea. One of the most recognizable constellations is the Big Dipper, which is part of the constellation Ursa Major. Constellations also serve to help describe the location of other objects in the sky.
| Physical sciences | Basics | Astronomy |
2776748 | https://en.wikipedia.org/wiki/Spinning%20mule | Spinning mule | The spinning mule is a machine used to spin cotton and other fibres. They were used extensively from the late 18th to the early 20th century in the mills of Lancashire and elsewhere. Mules were worked in pairs by a minder, with the help of two boys: the little piecer and the big or side piecer. The carriage carried up to 1,320 spindles and could be long, and would move forward and back a distance of four times a minute.
It was invented between 1775 and 1779 by Samuel Crompton. The self-acting (automatic) mule was patented by Richard Roberts in 1825. At its peak, there were 5,000,000 mule spindles in Lancashire alone. Modern versions are still in production and are used to spin woollen yarns from noble fibres such as cashmere, ultra-fine merino and alpaca for the knitted textile market.
The spinning mule spins textile fibres into yarn by an intermittent process. In the draw stroke, the roving is pulled through rollers and twisted; on the return it is wrapped onto the spindle. Its rival, the throstle frame or ring frame, uses a continuous process, where the roving is drawn, twisted and wrapped in one action. The mule was the most common spinning machine from 1790 until about 1900 and was still used for fine yarns until the early 1980s. In 1890, a typical cotton mill would have over 60 mules, each with 1,320 spindles, which would operate four times a minute for 56 hours a week.
History
Before the 1770s, textile production was a cottage industry using flax and wool. Weaving was a family activity. The children and women would card the fibre – break up and clean the disorganised fluff into long bundles. The women would then spin these rough rovings into yarn wound onto a spindle. The male weaver would use a frame loom to weave this into cloth. This was then tentered in the sun to bleach it. The invention by John Kay of the flying shuttle made the loom twice as productive, causing the demand for cotton yarn to vastly exceed what traditional spinners could supply.
There were two types of spinning wheel: the simple wheel, which uses an intermittent process, and the more refined Saxony wheel, which drives a differential spindle and flyer with a heck (an apparatus that guides the thread to the reels) in a continuous process. These two wheels became the starting point of technological development. Businessmen such as Richard Arkwright employed inventors to find solutions that would increase the amount of yarn spun, then took out the relevant patents.
The spinning jenny allowed a group of eight spindles to be operated together. It mirrored the simple wheel; the rovings were clamped, and a frame moved forward stretching and thinning the roving. A wheel was rapidly turned as the frame was pushed back, and the spindles rotated, twisting the rovings into yarn and collecting it on the spindles. The spinning jenny was effective and could be operated by hand, but it produced weaker thread that could be used only for the weft part of the cloth. (Because the side-to-side weft does not have to be stretched on a loom in the way that the warp is, it can generally be less strong.)
The throstle and the later water frame pulled the rovings through a set of attenuating rollers. Spinning at differing speeds, these pulled the thread continuously while other parts twisted it as it wound onto the heavy spindles. This produced thread suitable for warp, but the multiple rollers required much more energy input and demanded that the device be driven by a water wheel. The early water frame, however, had only a single spindle. Combining ideas from these two system inspired the spinning mule.
The increased supply of muslin inspired developments in loom design such as Edmund Cartwright's power loom. Some spinners and handloom weavers opposed the perceived threat to their livelihood: there were frame-breaking riots and, in 1811–13, the Luddite riots. The preparatory and associated tasks allowed many children to be employed until this was regulated.
Development over the next century and a half led to an automatic mule and to finer and stronger yarn. The ring frame, originating in New England in the 1820s, was little used in Lancashire until the 1890s. It required more energy and could not produce the finest counts.
The first mule
Samuel Crompton invented the spinning mule in 1779, so called because it is a hybrid of Arkwright's water frame and James Hargreaves's spinning jenny in the same way that a mule is the product of crossbreeding a female horse with a male donkey. The spinning mule has a fixed frame with a creel of cylindrical bobbins to hold the roving, connected through the headstock to a parallel carriage with the spindles. On the outward motion, the rovings are paid out through attenuating rollers and twisted. On the return, the roving is clamped and the spindles are reversed to take up the newly spun thread.
Crompton built his mule from wood. Although he used Hargreaves' ideas of spinning multiple threads and of attenuating the roving with rollers, it was he who put the spindles on the carriage and fixed a creel of roving bobbins on the frame. Both the rollers and the outward motion of the carriage remove irregularities from the rove before it is wound on the spindle. When Arkwright's patents expired, the mule was developed by several manufacturers.
Crompton's first mule had 48 spindles and could produce of 60s thread a day. This demanded a spindle speed of 1,700 rpm, and a power input of .
The mule produced strong, thin yarn, suitable for any kind of textile, warp or weft. It was first used to spin cotton, then other fibres.
Samuel Crompton could not afford to patent his invention. He sold the rights to David Dale and returned to weaving. Dale patented the mule and profited from it.
Improvements
Crompton's machine was largely built of wood, using bands and pulleys for the driving motions. After his machine was public, he had little to do with its development. Henry Stones, a mechanic from Horwich, constructed a mule using toothed gearing and, importantly, metal rollers. Baker of Bury worked on drums, and Hargreaves used parallel scrolling to achieve smoother acceleration and deceleration.
In 1790, William Kelly of Glasgow used a new method to assist the draw stroke. First animals, and then water, was used as the prime mover. Wright of Manchester moved the headstock to the centre of the machine, allowing twice as many spindles; a squaring band was added to ensure the spindles came out in a straight line. He was in conversation with John Kennedy about the possibility of a self-acting mule. Kennedy, a partner in McConnell & Kennedy machine makers in Ancoats, was concerned with building ever larger mules. McConnell & Kennedy ventured into spinning when they were left with two unpaid-for mules; their firm prospered and eventually merged into the Fine Spinners & Doublers Association. In 1793, John Kennedy addressed the problem of fine counts. With these counts, the spindles on the return traverse needed to rotate faster than on the outward traverse. He attached gears and a clutch to implement this motion.
William Eaton, in 1818, improved the winding of the thread by using two faller wires and performing a backing off at the end of the outward traverse. All these mules had been worked by the strength of the operatives. The next improvement was a fully automatic mule.
Roberts' self-acting mule
Richard Roberts took out his first patent in 1825 and a second in 1830. The task he had set himself was to design a self-actor, a self-acting or automatic spinning mule. Roberts is also known for the Roberts Loom, which was widely adopted because of its reliability. The mule in 1820 still needed manual assistance to spin a consistent thread; a self-acting mule would need:
A reversing mechanism that would unwind a spiral of yarn on the top of each spindle, before commencing the winding of a new stretch
A faller wire that would ensure the yarn was wound into a predefined form such as a cop
An appliance to vary the speed of revolution of the spindle, in accordance with the diameter of thread on that spindle
A counter faller under the thread was made to rise to take in the slack caused by backing off. This could be used with the top faller wire to guide the yarn to the correct place on the cop. These were controlled by levers and cams and an inclined plane called the shaper. The spindle speed was controlled by a drum and weighted ropes, as the headstock moved the ropes twisted the drum, which using a tooth wheel turned the spindles. None of this would have been possible using the technology of Crompton's time, fifty years earlier.
With the invention of the self actor, the hand-operated mule was increasingly referred to as a mule-jenny.
Oldham counts
Oldham counts refers to the medium thickness cotton that was used for general purpose cloth. Roberts did not profit from his self-acting spinning mule, but on the expiry of the patent other firms took forward the development, and the mule was adapted for the counts it spun. Initially Roberts' self-actor was used for coarse counts (Oldham Counts), but the mule-jenny continued to be used for the very finest counts (Bolton counts) until the 1890s and beyond.
Bolton counts
Bolton specialised in fine count cotton, and its mules ran more slowly to put in the extra twist. The mule jenny allowed for this gentler action but in the 20th century additional mechanisms were added to make the motion more gentle, leading to mules that used two or even three driving speeds. Fine counts needed a softer action on the winding, and relied on manual adjustment to wind the chase or top of the perfect cop.
Woollen mules
Spinning wool is a different process as the variable lengths of the individual fibres means that they are unsuitable for attenuation by roller drafting. For this reason, woolen fibres are carded using condenser cards which rub the carded fibres together rather than drafting them. They are then spun on mule-type machines which have no roller drafting, but create the draft by the spindles receding from the delivery rollers whilst that latter, having paid out a short length of roving, are held stationary. Such mules are often complex involving multiple spindles speeds, receding motions, etc. to ensure optimum treatment of the yarn.
Condenser spinning
Condenser spinning was developed to enable the short fibres produced as waste from the combing of fine cottons, to be spun into a soft, coarse yarns suitable for sheeting, blankets etc. Only approximately 2% of the mule spindles in Lancashire were Condenser spindles, but many more Condenser mules survive today as these were the last spindles regularly at work., and the mules are similar. Helmshore Mills was a cotton waste mule spinning mill.
Current usage
Mules are still in use for spinning woolen and alpaca, and being produced across the world. In Italy for example by Bigagli and Cormatex
Operation of a mule
Watch video demonstration #1
Mule spindles rest on a carriage that travels on a track a distance of , while drawing out and spinning the yarn. On the return trip, known as putting up, as the carriage moves back to its original position, the newly spun yarn is wound onto the spindle in the form of a cone-shaped cop. As the mule spindle travels on its carriage, the roving which it spins is fed to it through rollers geared to revolve at different speeds to draw out the yarn.
Marsden in 1885 described the processes of setting up and operating a mule. Here is his description, edited slightly.
The creel holds bobbins containing rovings. The rovings are
passed through small guide-wires, and between the three pairs of drawing-rollers.
The first pair takes hold of the roving, to draw the roving or sliver from the bobbin, and deliver it to the next pair.
The motion of the middle pair is slightly quicker than the first, but only sufficiently so to keep the roving uniformly tense
The front pair, running much more quickly, draws out (attenuates) the roving so it is equal throughout.
Connection is then established between the attenuated rovings and the spindles. When the latter are bare, as in a new mule, the spindle-driving motion is put into gear, and the attendants wind upon each spindle a short length of yarn from a cop held in the hand. The drawing-roller motion is placed in gear, and the rollers soon present lengths of attenuated roving. These are attached to the threads on the spindles, by simply placing the threads in contact with the un-twisted roving. The different parts of the machine are next simultaneously started, when the whole works in harmony together.
The back rollers pull the sliver from the bobbins, and passing it to the succeeding pairs, whose differential speeds attenuate it to the required degree of fineness. As it is delivered in front, the spindles, revolving at a rate of 6,000–9,000 rpm twist the hitherto loose fibres together, thus forming a thread.
Whilst this is going on, the spindle carriage is being drawn away from the rollers, at a pace very slightly exceeding the rate at which the roving is coming forth. This is called the gain of the carriage, its purpose being to eliminate all irregularities in the fineness of the thread. Should a thick place in the roving come through the rollers, it would resist the efforts of the spindle to twist it; and, if passed in this condition, it would seriously deteriorate the quality of the yarn, and impede subsequent operations. As, however, the twist, spreading itself over the level thread, gives firmness to this portion, the thick and untwisted part yields to the draught of the spindle, and, as it approaches the tenuity of the remainder, it receives the twist it had hitherto refused to take. The carriage, which is borne upon wheels, continues its outward progress, until it reaches the extremity of its traverse, which is from the roller beam. The revolution of the spindles cease, the drawing rollers stop.
Backing-off commences. This process is the unwinding of the several turns of the yarn, extending from the top of the cop in process of formation to the summit of the spindle. As this proceeds, the faller- wire, which is placed over and guides the threads upon the cop, is depressed; the counter-faller at the same time rising, the slack unwound from the spindles is taken up, and the threads are prevented from running into snarls. Backing-off is completed.
The carriage commences to run inwards; that is, towards the rollerbeam. This is called putting up. The spindles wind on the yarn at a uniform rate. The speed of revolution of the spindle must vary, as the faller is guiding the thread upon the larger or smaller
diameter of the cone of the cop. Immediately the winding is finished, the depressed faller rises, the counter-faller is put down.
These movements are repeated until the cops on each spindle are perfectly formed: the set is completed. A stop-motion paralyses every action of the machine, rendering it necessary to doff or strip the spindles, and to commence anew.
Doffing is performed by the piercers thrutching, that is raising, the cops partially up the spindles, whilst the carriage is out. The minder then depressing the faller, so far as to guide the threads upon the bare spindle below. A few turns are wound onto the spindle, to fix the threads to the bare spindles for a new set. The cops are removed and collected into cans or baskets, and subsequently delivered to the warehouse. The remainder of the "draw" or "stretch," as the length of spun yarn is called when the carriage is out, is then wound upon the spindles as the carriage is run up to the roller beam. Work then commences anew.
The doffing took only a few minutes, the piecers would run the length of the mule gate thrutching five spindles a time, and the doffing involved lifting four cops from the spindles with the right hand and piling them on the left forearm and hand. To get a firm cop bottom, the minder would whip the first few layers of yarn. After the first few draws the minder would stop the mule at the start of an inward run and take it in slowly depressing and releasing the faller wire several times. Alternatively, a starch paste could be skilfully applied to the first few layers of yarn by the piecers – and later a small paper tube was dropped over spindle – this slowed down the doffing operation and extra payment was negotiated by the minders.
Duties of the operatives
A pair of mules would be manned by a person called the minder and two boys called the side piecer and the little piecer. They worked barefoot in humid temperatures; the minder and the little piecer worked the minder half of the mule. The minder would make minor adjustments to his mules to the extent that each mule worked differently. They were specialists in spinning, and were answerable only to the gaffer and under-gaffer who were in charge of the floor and with it the quantity and quality of the yarn that was produced. Bobbins of rovings came from the carder in the blowing room delivered by a bobbin carrier who was part of the carder's staff, and yarn was hoisted down to the warehouse by the warehouseman's staff. Delineation of jobs was rigid and communication would be through the means of coloured slips of paper written on in indelible pencil.
Creeling involved replacing the rovings bobbins in a section of the mule without stopping the mule. On very coarse counts a bobbin lasted two days but on fine count it could last for 3 weeks. To creel, the creeler stood behind the mule, placing new bobbins on the shelf above the creel. As the bobbin ran empty he would pick it off its skewer in the creel unreeling 30 cm or so of roving, and drop it into a skip. With his left hand, he would place on the new bobbin onto the skewer from above and with his right hand twist in the new roving into the tail of the last.
Piecing involved repairing sporadic yarn breakages. At the rollers, the broken yarn would be caught on the underclearer (or fluker rod on Bolton mules), while at the spindle it would knot itself into a whorl on the spindle tip. If the break happened on the winding stroke the spindle might have to be stopped while the thread was found. The number of yarn breakages was dependent on the quality of the roving, and quality cotton led to fewer breakages. Typical 1,200 spindle mules of the 1920s would experience 5 to 6 breakages a minute. The two piecers would thus need to repair the thread within 15 to 20 seconds while the mule was in motion but once they had the thread it took under three seconds. The repair actually involved a slight rolling of the forefinger against the thumb.
Doffing has already been described.
Cleaning was important and until a formal ritual had been devised it was a dangerous operation. The vibration in a mule threw a lot of short fibres (or fly) into the air. It tended to accumulate on the carriage behind the spindles and in the region of the drafting rollers. Piking the stick meant placing the hand though the yarnsheet, and unclipping two sticks of underclearer rollers from beneath the drafting rollers, drawing them through the gap between two ends, stripping them of fly and replacing them on the next inward run. Cleaning the carriage top was far more dangerous. The minder would stop the mule on the outward run, and raise his hands above his head. The piecers would enter under the yarn sheet with a scavenger cloth on the carriage spindle rail and a brush on the roller beam, and run bent double the entire length of the mule, avoiding the rails and draw bands, and not letting themselves touch the yarn sheet. When they had finished they would run to agreed positions of safety where the minder could see both of them, and the minder would unclip the stang and start the mule. Before this ritual was devised, boys had been crushed. The mule was long, the minder's eyesight might not have been good, the air in the mill was clouded with fly and another minder's boys might have been mistaken for his. The ritual became encoded in law.
Key components
Drawing rollers
Faller and counter faller
Quadrant
Terminology
Social and economic impact
The spinning inventions were significant in enabling a great expansion to occur in the production of textiles, particularly cotton ones. Cotton and iron were leading sectors in the Industrial Revolution. Both industries underwent a great expansion at about the same time, which can be used to identify the start of the Industrial Revolution.
The 1790 mule was operated by brute force: the spinner drawing and pushing the frame while attending to each spindle. Home spinning was the occupation of women and girls, but the strength needed to operate a mule caused it to be the activity of men. Hand loom weaving, however, had been a man's occupation but in the mill it could and was done by girls and women. Spinners were the bare-foot aristocrats of the factory system. It replaced decentralised cottage industries with centralised factory jobs, driving economic upheaval and urbanisation.
Mule spinners were the leaders in unionism within the cotton industry; the pressure to develop the self-actor or self-acting mule was partly to open the trade to women. It was in 1870 that the first national union was formed.
The wool industry was divided into woollen and worsted. It lagged behind cotton in adopting new technology. Worsted tended to adopt Arkwright water frames which could be operated by young girls, and woollen adopted the mule.
Mule-spinners' cancer
Circa 1900 there was a high incidence of scrotal cancer detected in former mule spinners. It was limited to cotton mule spinners and did not affect woollen or condenser mule spinners. The cause was attributed to the blend of vegetable and mineral oils used to lubricate the spindles. The spindles, when running, threw out a mist of oil at crotch height, that was captured by the clothing of anyone piecing an end. In the 1920s much attention was given to this problem. Mules had used this mixture since the 1880s, and cotton mules ran faster and hotter than the other mules, and needed more frequent oiling. The solution was to make it a statutory requirement to use only vegetable oil or white mineral oils, which were believed to be non-carcinogenic. By then cotton mules had been superseded by the ring frame and the industry was contracting, so it was never established whether these measures were effective.
| Technology | Spinning | null |
2776778 | https://en.wikipedia.org/wiki/Barbary%20sheep | Barbary sheep | The Barbary sheep (Ammotragus lervia), also known as aoudad (pronounced [ˈɑʊdæd]), is a species of caprine native to rocky mountains in North Africa and parts of West Africa. While this is the only species in genus Ammotragus, six subspecies have been described. Although it is rare in its native North Africa, it has been introduced to North America, southern Europe, and elsewhere. It is also known in the Berber language as waddan or arwi, and in former French territories as the mouflon.
Description
Barbary sheep stand tall at the shoulder, with a length around , and weigh . They are sandy-brown, darkening with age, with a slightly lighter underbelly and a darker line along the back. Upper parts and the outer parts of the legs are a uniform reddish- or grayish-brown. Some shaggy hair is on the throat (extending down to the chest in males) with a sparse mane. Their horns have a triangular cross-section. The horns curve outward, backward, then inward, and can exceed in length. The horns are fairly smooth, with slight wrinkles evident at the base as the animal matures.
Range
Natural range
Barbary sheep are endemic to regions of Northern Africa primarily surrounding the barren center of the Sahara Desert. Countries and territories where aoudad may be found include Algeria, Chad (north), Egypt, Libya, Mali (north), Mauritania, Morocco, Niger, Tunisia and Western Sahara. West of the Nile, they can be found in Sudan; east of the Nile, in the Red Sea Hills. The now-extinct Ancient Egyptian corkscrew-horned sheep (Ovis longipes palaeoaegyptiacus) was also thought to be a subspecies of wild barbary sheep.
Populations within its native range have been decreasing due to hunting, legal and otherwise, and destruction of habitat.
Introduced populations
Barbary sheep have been introduced to southeastern Spain and the southwestern United States.
They have become common in a limited region of southeastern Spain, since its introduction in 1970 to Sierra Espuña Regional Park as a game species. Its adaptability enabled it to colonize nearby areas quickly, and private game estates provided other centers of dispersion. The species is currently expanding, according to recent field surveys, now being found in the provinces of Alicante, Almería, Granada, and Murcia. The species is a potential competitor to native ungulates inhabiting the Iberian Peninsula, and has also been introduced to La Palma (in the Canary Islands), and has spread throughout the northern and central parts of the island, where it is a serious threat to endemic vegetation. The aoudad has also been introduced in Croatia several times, where there is a population in Mosor.
Although the species has not yet been recorded in Australia, it is considered a pest species in Queensland with the potential to establish in the wild.
Taxonomy
A. lervia is the only species in the genus Ammotragus. However, some authors include this genus in the goat genus Capra, together with the sheep genus Ovis.
The subspecies are found allopatrically in various parts of North Africa:
A. l. lervia Pallas, 1777 (vulnerable)
A. l. ornata I. Geoffroy Saint-Hilaire, 1827 (Egyptian Barbary sheep, thought to be extinct in the wild but still found in the eastern desert of Egypt)
A. l. sahariensis Rothschild, 1913 (vulnerable)
A. l. blainei Rothschild, 1913 (vulnerable)
A. l. angusi Rothschild, 1921 (vulnerable)
A. l. fassini Lepri, 1930 (vulnerable)
Habitats
Barbary sheep are found in arid mountainous areas where they graze and browse grasses, bushes, and lichens. They are able to obtain all their metabolic water from food, but if liquid water is available, they drink and wallow in it. Barbary sheep are crepuscular - active in the early morning and late afternoon and rest in the heat of the day. They are very agile and can achieve a standing jump over . They are well adapted to their habitat, which consist of steep, rocky mountains and canyons. They often flee at the first sign of danger, typically running uphill. They are extremely nomadic and travel constantly via mountain ranges. Their main predators in North Africa were the Barbary leopard, Barbary lion, and caracal, but now humans, feral dogs, competition due to overgrazing by domestic animals and drought threaten their populations.
Names
The binomial name Ammotragus lervia derives from the Greek ἄμμος ámmos ("sand", referring to the sand-coloured coat) and τράγος trágos ("goat").
Lervia derives from the wild sheep of northern Africa described as "lerwee" by Rev. T. Shaw in his "Travels and Observations" about parts of Barbary and Levant.
The Spanish named this sheep the arruis, from Berber arrwis, and the Spanish Legion even used it as a mascot for a time.
Aoudad () is the name for this sheep used by the Berbers, a North African people, and it is also called arui and waddan (in Libya).
Gallery
| Biology and health sciences | Bovidae | Animals |
34311891 | https://en.wikipedia.org/wiki/Cephalopod%20beak | Cephalopod beak | All extant cephalopods have a two-part beak, or rostrum, situated in the buccal mass and surrounded by the muscular head appendages. The dorsal (upper) mandible fits into the ventral (lower) mandible and together they function in a scissor-like fashion. The beak may also be referred to as the mandibles or jaws. These beaks are different from bird beaks because they crush bone while most bird beaks do not.
Fossilised remains of beaks are known from a number of cephalopod groups, both extant and extinct, including squids, octopuses, belemnites, and vampyromorphs. Aptychi – paired plate-like structures found in ammonites – may also have been jaw elements.
Composition
Composed primarily of chitin and cross-linked proteins, beaks are more-or-less indigestible and are often the only identifiable cephalopod remains found in the stomachs of predatory species such as sperm whales. Cephalopod beaks gradually become less stiff as one moves from the tip to the base, a gradient that results from differing chemical composition. In hydrated beaks of the Humboldt squid (Dosidicus gigas) this stiffness gradient spans two orders of magnitude.
Measurements
The abbreviations LRL and URL are commonly used in teuthology to refer to lower rostral length and upper rostral length, respectively. These are the standard measures of beak size in Decapodiformes; hood length is preferred for Octopodiformes. They can be used to estimate the mantle length and total body weight of the original animal as well as the total ingested biomass of the species.
| Biology and health sciences | Gastrointestinal tract | Biology |
34317494 | https://en.wikipedia.org/wiki/QED%20vacuum | QED vacuum | The QED vacuum or quantum electrodynamic vacuum is the field-theoretic vacuum of quantum electrodynamics. It is the lowest energy state (the ground state) of the electromagnetic field when the fields are quantized. When the Planck constant is hypothetically allowed to approach zero, QED vacuum is converted to classical vacuum, which is to say, the vacuum of classical electromagnetism.
Another field-theoretic vacuum is the QCD vacuum of the Standard Model.
Fluctuations
The QED vacuum is subject to fluctuations about a dormant zero average-field condition; Here is a description of the quantum vacuum:
Virtual particles
It is sometimes attempted to provide an intuitive picture of virtual particles based upon the Heisenberg energy-time uncertainty principle:
(where and are energy and time variations, and the Planck constant divided by 2) arguing along the lines that the short lifetime of virtual particles allows the "borrowing" of large energies from the vacuum and thus permits particle generation for short times.
This interpretation of the energy-time uncertainty relation is not universally accepted, however. One issue is the use of an uncertainty relation limiting measurement accuracy as though a time uncertainty determines a "budget" for borrowing energy . Another issue is the meaning of "time" in this relation, because energy and time (unlike position and momentum , for example) do not satisfy a canonical commutation relation (such as ). Various schemes have been advanced to construct an observable that has some kind of time interpretation, and yet does satisfy a canonical commutation relation with energy. The many approaches to the energy-time uncertainty principle are a continuing subject of study.
Quantization of the fields
The Heisenberg uncertainty principle does not allow a particle to exist in a state in which the particle is simultaneously at a fixed location, say the origin of coordinates, and has also zero momentum. Instead the particle has a range of momentum and spread in location attributable to quantum fluctuations; if confined, it has a zero-point energy.
An uncertainty principle applies to all quantum mechanical operators that do not commute. In particular, it applies also to the electromagnetic field. A digression follows to flesh out the role of commutators for the electromagnetic field.
The standard approach to the quantization of the electromagnetic field begins by introducing a vector potential and a scalar potential to represent the basic electromagnetic electric field and magnetic field using the relations: The vector potential is not completely determined by these relations, leaving open a so-called gauge freedom. Resolving this ambiguity using the Coulomb gauge leads to a description of the electromagnetic fields in the absence of charges in terms of the vector potential and the momentum field , given by: where is the electric constant of the SI units. Quantization is achieved by insisting that the momentum field and the vector potential do not commute. That is, the equal-time commutator is: where , are spatial locations, is the reduced Planck constant, is the Kronecker delta and is the Dirac delta function. The notation denotes the commutator.
Quantization can be achieved without introducing the vector potential, in terms of the underlying fields themselves: where the circumflex denotes a Schrödinger time-independent field operator, and is the antisymmetric Levi-Civita tensor.
Because of the non-commutation of field variables, the variances of the fields cannot be zero, although their averages are zero. The electromagnetic field has therefore a zero-point energy, and a lowest quantum state. The interaction of an excited atom with this lowest quantum state of the electromagnetic field is what leads to spontaneous emission, the transition of an excited atom to a state of lower energy by emission of a photon even when no external perturbation of the atom is present.
Electromagnetic properties
As a result of quantization, the quantum electrodynamic vacuum can be considered as a material medium. It is capable of vacuum polarization. In particular, the force law between charged particles is affected. The electrical permittivity of quantum electrodynamic vacuum can be calculated, and it differs slightly from the simple of the classical vacuum. Likewise, its permeability can be calculated and differs slightly from . This medium is a dielectric with relative dielectric constant > 1, and is diamagnetic, with relative magnetic permeability < 1. Under some extreme circumstances in which the field exceeds the Schwinger limit (for example, in the very high fields found in the exterior regions of pulsars), the quantum electrodynamic vacuum is thought to exhibit nonlinearity in the fields. Calculations also indicate birefringence and dichroism at high fields. Many of electromagnetic effects of the vacuum are small, and only recently have experiments been designed to enable the observation of nonlinear effects. PVLAS and other teams are working towards the needed sensitivity to detect QED effects.
Attainability
A perfect vacuum is itself only attainable in principle. It is an idealization, like absolute zero for temperature, that can be approached, but never actually realized:
Virtual particles make a perfect vacuum unrealizable, but leave open the question of attainability of a quantum electrodynamic vacuum or QED vacuum. Predictions of QED vacuum such as spontaneous emission, the Casimir effect and the Lamb shift have been experimentally verified, suggesting QED vacuum is a good model for a high quality realizable vacuum. There are competing theoretical models for vacuum, however. For example, quantum chromodynamic vacuum includes many virtual particles not treated in quantum electrodynamics. The vacuum of quantum gravity treats gravitational effects not included in the Standard Model. It remains an open question whether further refinements in experimental technique ultimately will support another model for realizable vacuum.
| Physical sciences | Quantum mechanics | Physics |
5063242 | https://en.wikipedia.org/wiki/Potassium%20acetate | Potassium acetate | Potassium acetate (also called potassium ethanoate), (CH3COOK) is the potassium salt of acetic acid. It is a hygroscopic solid at room temperature.
Preparation
It can be prepared by treating a potassium-containing base such as potassium hydroxide or potassium carbonate with acetic acid:
CH3COOH + KOH → CH3COOK + H2O
This sort of reaction is known as an acid-base neutralization reaction.
At saturation, the sesquihydrate in water solution (CH3COOK·1½H2O) begins to form semihydrate at 41.3 °C.
Applications
Deicing
Potassium acetate (as a substitute for calcium chloride or magnesium chloride) can be used as a deicer to remove ice or prevent its formation. It offers the advantage of being less aggressive on soils and much less corrosive: for this reason, it is preferred for airport runways although it is more expensive.
Fire extinguishing
Potassium acetate is the extinguishing agent used in Class K fire extinguishers because of its ability to cool and form a crust over burning oils.
Food additive
Potassium acetate is used in processed foods as a preservative and acidity regulator. In the European Union, it is labeled by the E number E261; it is also approved for usage in the USA, Australia, and New Zealand. Potassium hydrogen diacetate (CAS #) with formula KH(OOCCH3)2 is a related food additive with the same E number as potassium acetate.
Medicine and biochemistry
In medicine, potassium acetate is used as part of electrolyte replacement protocols in the treatment of diabetic ketoacidosis because of its ability to break down to bicarbonate to help neutralize the acidotic state.
In molecular biology, potassium acetate is used to precipitate Sodium dodecyl sulfate (SDS) and SDS-bound proteins to allow their removal from DNA.
Potassium acetate is used in mixtures applied for tissue preservation, fixation, and mummification. Most museums today use a formaldehyde-based method recommended by Kaiserling in 1897 which contains potassium acetate. This process was used to soak Lenin's corpse.
Use in executions
Potassium acetate was incorrectly used in place of potassium chloride when putting a prisoner to death in Oklahoma in January 2015. Charles Frederick Warner was executed on January 15, 2015 with potassium acetate; this was not public knowledge until the scheduled execution of Richard Glossip was called off. In August 2017, the U.S. state of Florida executed Mark James Asay using a combination of etomidate, rocuronium bromide, and potassium acetate. The drug was also used in the February 2023 execution of Donald Dillbeck, once again in combination with etomidate and rocuronium bromide.
Industry
Potassium acetate is used as a catalyst in the production of polyurethanes.
Historical
It is used as a diuretic and urinary alkalizer. Before modern chemistry, it was variously called terra foliata tartari, sal Sennerti, tartarus regeneratus, arcanum tartari and sal diureticus. In 1760 it was used in the preparation of Cadet's fuming liquid ((CH3)2As)2O, the first organometallic compound ever produced.
| Physical sciences | Acetates | Chemistry |
5064554 | https://en.wikipedia.org/wiki/Lead%28II%29%20chromate | Lead(II) chromate | Lead(II) chromate is an inorganic compound with the chemical formula . It is a bright yellow solid that is very poorly soluble in water. It occurs also as the mineral crocoite. It is used as a pigment.
Structure
Two polymorphs of lead chromate are known, orthorhombic and the more stable monoclinic form. Monoclinic lead chromate is used in paints under the name chrome yellow, and many other names. Lead chromate adopts the monazite structure, meaning that the connectivity of the atoms is very similar to other compounds of the type . Pb(II) has a distorted coordination sphere being surrounded by eight oxides with Pb-O distances ranging from 2.53 to 2.80 Å. The chromate anion is tetrahedral, as usual. Unstable polymorphs of lead chromate are the greenish yellow orthorhombic form and a red-orange tetragonal form.
Applications
Approximately 37,000 tons were produced in 1996. The main applications are as a pigment in paints, under the name chrome yellow.
Preparation
Lead(II) chromate can be produced by treating sodium chromate with lead salts such as lead(II) nitrate or by combining lead(II) oxide with chromic acid.
Related lead sulfochromate pigments are produced by the replacement of some chromate by sulfate, resulting in a mixed lead-chromate-sulfate compositions . This replacement is possible because sulfate and chromate are isostructural. Since sulfate is colorless, sulfochromates with high values of x are less intensely colored than lead chromate. In some cases, chromate is replaced by molybdate.
Reactions
Heating in hydroxide solution produces chrome red, a red or orange powder made by PbO and . Also, in hydroxide solution lead chromate slowly dissolves forming plumbite complex.
Safety hazards
Despite containing both lead and hexavalent chromium, lead chromate is not acutely lethal because of its very low solubility. The LD50 for rats is only 5,000 mg/kg. Lead chromate must be treated with great care in its manufacture, the main concerns being dust of the chromate precursor. Lead chromate is highly regulated in advanced countries. As one of the greatest threats comes from inhalation of particles, so much effort has been devoted to production of low-dust forms of the pigment.
In the 1800s, the product was used to impart a bright yellow color to some types of candy. It is used (illegally) to enhance the color of certain spices, particularly turmeric, particularly in Bangladesh.
Unlike other lead-based paint pigments, lead chromate is still widely used, especially in road marking paint.
In 2023 and 2024, consumption of adulterated cinnamon led to at least 136 cases of lead toxicity in children in the United States as reported by the US Centers for Disease Control and Prevention. The affected products were recalled. The US Food and Drug Administration determined that the ratio of lead to chromium in the cinammon indicated that lead chromate had been added to the cinnamon.
| Physical sciences | Metallic oxyanions | Chemistry |
21242393 | https://en.wikipedia.org/wiki/Musa%20balbisiana | Musa balbisiana | Musa balbisiana, also known simply as plantain, is a wild-type species of banana. It is one of the ancestors of modern cultivated bananas, along with Musa acuminata.
Description
It grows lush leaves in clumps with a more upright habit than most cultivated bananas. Flowers grow in inflorescences coloured red to maroon. The fruit are between blue and green. They are considered inedible because of the seeds they contain.
Taxonomy
It was first scientifically described in 1820 by the Italian botanist Luigi Aloysius Colla.
Distribution
It is native to eastern South Asia, the eastern regions of the Indian subcontinent, northern Southeast Asia, and southern China. Introduced populations exist in the wild, far outside its native range.
Uses
It is assumed that wild bananas were cooked and eaten, as farmers would not have developed the cultivated banana otherwise. Seeded Musa balbisiana fruit are called butuhan ('with seeds') in the Philippines, and kluai tani (กล้วยตานี) in Thailand, where its leaves are used for packaging and crafts. Natural parthenocarpic clones occur through polyploidy and produce edible bananas, examples of which are wild saba bananas.
Genome
Musa balbisiana contributed the B genome to the cultivated banana.
Wang et al., 2019 provides a genome, evolutionary analysis and functional genomics analysis. Wang et al. find evolution increasing ethylene production in the domesticated form.
| Biology and health sciences | Tropical and tropical-like fruit | Plants |
21244265 | https://en.wikipedia.org/wiki/Sense%20of%20smell | Sense of smell | The sense of smell, or olfaction, is the special sense through which smells (or odors) are perceived. The sense of smell has many functions, including detecting desirable foods, hazards, and pheromones, and plays a role in taste.
In humans, it occurs when an odor binds to a receptor within the nasal cavity, transmitting a signal through the olfactory system. Glomeruli aggregate signals from these receptors and transmit them to the olfactory bulb, where the sensory input will start to interact with parts of the brain responsible for smell identification, memory, and emotion.
There are many different things which can interfere with a normal sense of smell, including damage to the nose or smell receptors, anosmia, upper respiratory infections, traumatic brain injury, and neurodegenerative disease.
History of study
Early scientific study of the sense of smell includes the extensive doctoral dissertation of Eleanor Gamble, published in 1898, which compared olfactory to other stimulus modalities, and implied that smell had a lower intensity discrimination.
As the Epicurean and atomistic Roman philosopher Lucretius (1stcentury BCE) speculated, different odors are attributed to different shapes and sizes of "atoms" (odor molecules in the modern understanding) that stimulate the olfactory organ.
A modern demonstration of that theory was the cloning of olfactory receptor proteins by Linda B. Buck and Richard Axel (who were awarded the Nobel Prize in 2004), and subsequent pairing of odor molecules to specific receptor proteins. Each odor receptor molecule recognizes only a particular molecular feature or class of odor molecules. Mammals have about a thousand genes that code for odor reception. Of the genes that code for odor receptors, only a portion are functional. Humans have far fewer active odor receptor genes than other primates and other mammals. In mammals, each olfactory receptor neuron expresses only one functional odor receptor. Odor receptor nerve cells function like a key–lock system: if the airborne molecules of a certain chemical can fit into the lock, the nerve cell will respond.
There are, at present, a number of competing theories regarding the mechanism of odor coding and perception. According to the shape theory, each receptor detects a feature of the odor molecule. The weak-shape theory, known as the odotope theory, suggests that different receptors detect only small pieces of molecules, and these minimal inputs are combined to form a larger olfactory perception (similar to the way visual perception is built up of smaller, information-poor sensations, combined and refined to create a detailed overall perception).
According to a new study, researchers have found that a functional relationship exists between molecular volume of odorants and the olfactory neural response. An alternative theory, the vibration theory proposed by Luca Turin, posits that odor receptors detect the frequencies of vibrations of odor molecules in the infrared range by quantum tunnelling. However, the behavioral predictions of this theory have been called into question. There is no theory yet that explains olfactory perception completely.
Function
Taste
Flavor perception is an aggregation of auditory, taste, haptic, and smell sensory information. Retronasal smell plays the biggest role in the sensation of flavor. During the process of mastication, the tongue manipulates food to release odorants. These odorants enter the nasal cavity during exhalation. The smell of food has the sensation of being in the mouth because of co-activation of the motor cortex and olfactory epithelium during mastication.
Smell, taste, and trigeminal receptors (also called chemesthesis) together contribute to flavor. The human tongue can distinguish only among five distinct qualities of taste, while the nose can distinguish among hundreds of substances, even in minute quantities. It is during exhalation that the smell's contribution to flavor occurs, in contrast to that of proper smell, which occurs during the inhalation phase of breathing. The olfactory system is the only human sense that bypasses the thalamus and connects directly to the forebrain.
Hearing
Smell and sound information has been shown to converge in the olfactory tubercles of rodents. This neural convergence is proposed to give rise to a perception termed smound. Whereas a flavor results from interactions between smell and taste, a smound may result from interactions between smell and sound.
Inbreeding avoidance
The MHC genes (known as HLA in humans) are a group of genes present in many animals and important for the immune system; in general, offspring from parents with differing MHC genes have a stronger immune system. Fish, mice, and female humans are able to smell some aspect of the MHC genes of potential sex partners and prefer partners with MHC genes different from their own. However, some research suggests that taking hormonal contraception can alter women's preference for partners with dissimilar MHC genes, thus resulting in a greater likelihood to choose partners with relatively similar MHC genes to their own. Sexual orientation can also influence preference for different body odors, and some studies suggest that preference may be influenced by the putative pheromones AND and EST.
Humans can detect blood relatives from olfaction. Mothers can identify by body odor their biological children but not their stepchildren. Pre-adolescent children can olfactorily detect their full siblings but not half-siblings or step siblings, and this might explain incest avoidance and the Westermarck effect. Functional imaging shows that this olfactory kinship detection process involves the frontal-temporal junction, the insula, and the dorsomedial prefrontal cortex, but not the primary or secondary olfactory cortices, or the related piriform cortex or orbitofrontal cortex.
Since inbreeding is detrimental, it tends to be avoided. In the house mouse, the major urinary protein (MUP) gene cluster provides a highly polymorphic scent signal of genetic identity that appears to underlie kin recognition and inbreeding avoidance. Thus, there are fewer matings between mice sharing MUP haplotypes than would be expected if there were random mating.
Guiding movement
Some animals use scent trails to guide movement, for example social insects may lay down a trail to a food source, or a tracking dog may follow the scent of its target. A number of scent-tracking strategies have been studied in different species, including gradient search or chemotaxis, anemotaxis, klinotaxis, and tropotaxis. Their success is influenced by the turbulence of the air plume that is being followed.
Genetics
Different people smell different odors, and most of these differences are caused by genetic differences. Although odorant receptor genes make up one of the largest gene families in the human genome, only a handful of genes have been linked conclusively to particular smells. For instance, the odorant receptor OR5A1 and its genetic variants (alleles) are responsible for our ability (or failure) to smell β-ionone, a key aroma in foods and beverages. Similarly, the odorant receptor OR2J3 is associated with the ability to detect the "grassy" odor, cis-3-hexen-1-ol. The preference (or dislike) of cilantro (coriander) has been linked to the olfactory receptor OR6A2.
Variability amongst vertebrates
The importance and sensitivity of smell varies among different organisms; most mammals have a good sense of smell, whereas most birds do not, except the tubenoses (e.g., petrels and albatrosses), certain species of new world vultures, and the kiwis. Also, birds have hundreds of olfactory receptors. Although, recent analysis of the chemical composition of volatile organic compounds (VOCs) from king penguin feathers suggest that VOCs may provide olfactory cues, used by the penguins to locate their colony and recognize individuals. Among mammals, it is well developed in the carnivores and ungulates, which must always be aware of each other, and in those that smell for their food, such as moles. Having a strong sense of smell is referred to as macrosmatic in contrast to having a weak sense of smell which is referred to as microsmotic.
Figures suggesting greater or lesser sensitivity in various species reflect experimental findings from the reactions of animals exposed to aromas in known extreme dilutions. These are, therefore, based on perceptions by these animals, rather than mere nasal function. That is, the brain's smell-recognizing centers must react to the stimulus detected for the animal to be said to show a response to the smell in question. It is estimated that dogs, in general, have an olfactory sense approximately ten thousand to a hundred thousand times more acute than a human's. This does not mean they are overwhelmed by smells our noses can detect; rather, it means they can discern a molecular presence when it is in much greater dilution in the carrier, air.
Scenthounds as a group can smell one- to ten-million times more acutely than a human, and bloodhounds, which have the keenest sense of smell of any dogs, have noses ten- to one-hundred-million times more sensitive than a human's. They were bred for the specific purpose of tracking humans, and can detect a scent trail a few days old. The second-most-sensitive nose is possessed by the Basset Hound, which was bred to track and hunt rabbits and other small animals.
Grizzly bears have a sense of smell seven times stronger than that of the bloodhound, essential for locating food underground. Using their elongated claws, bears dig deep trenches in search of burrowing animals and nests as well as roots, bulbs, and insects. Bears can detect the scent of food from up to eighteen miles away; because of their immense size, they often scavenge new kills, driving away the predators (including packs of wolves and human hunters) in the process.
The sense of smell is less developed in the catarrhine primates, and nonexistent in cetaceans, which compensate with a well-developed sense of taste. In some strepsirrhines, such as the red-bellied lemur, scent glands occur atop the head. In many species, smell is highly tuned to pheromones; a male silkworm moth, for example, can sense a single molecule of bombykol.
Fish, too, have a well-developed sense of smell, even though they inhabit an aquatic environment. Salmon utilize their sense of smell to identify and return to their home stream waters. Catfish use their sense of smell to identify other individual catfish and to maintain a social hierarchy. Many fishes use the sense of smell to identify mating partners or to alert to the presence of food.
Human smell abilities
Although conventional wisdom and lay literature, based on impressionistic findings in the 1920s, have long presented human smell as capable of distinguishing between roughly 10,000 unique odors, recent research has suggested that the average individual is capable of distinguishing over one trillion unique odors. Researchers in the most recent study, which tested the psychophysical responses to combinations of over 128 unique odor molecules with combinations composed of up to 30 different component molecules, noted that this estimate is "conservative" and that some subjects of their research might be capable of deciphering between a thousand trillion odorants, adding that their worst performer could probably still distinguish between 80million scents. Authors of the study concluded, "This is far more than previous estimates of distinguishable olfactory stimuli. It demonstrates that the human olfactory system, with its hundreds of different olfactory receptors, far out performs the other senses in the number of physically different stimuli it can discriminate." However, it was also noted by the authors that the ability to distinguish between smells is not analogous to being able to consistently identify them, and that subjects were not typically capable of identifying individual odor stimulants from within the odors the researchers had prepared from multiple odor molecules. In November 2014 the study was strongly criticized by Caltech scientist Markus Meister, who wrote that the study's "extravagant claims are based on errors of mathematical logic." The logic of his paper has in turn been criticized by the authors of the original paper.
Physiological basis in vertebrates
Main olfactory system
In humans and other vertebrates, smells are sensed by olfactory sensory neurons in the olfactory epithelium. The olfactory epithelium is made up of at least six morphologically and biochemically different cell types. The proportion of olfactory epithelium compared to respiratory epithelium (not innervated, or supplied with nerves) gives an indication of the animal's olfactory sensitivity. Humans have about of olfactory epithelium, whereas some dogs have . A dog's olfactory epithelium is also considerably more densely innervated, with a hundred times more receptors per square centimeter. The sensory olfactory system integrates with other senses to form the perception of flavor. Often, land organisms will have separate olfaction systems for smell and taste (orthonasal smell and retronasal smell), but water-dwelling organisms usually have only one system.
Molecules of odorants passing through the superior nasal concha of the nasal passages dissolve in the mucus that lines the superior portion of the cavity and are detected by olfactory receptors on the dendrites of the olfactory sensory neurons. This may occur by diffusion or by the binding of the odorant to odorant-binding proteins. The mucus overlying the epithelium contains mucopolysaccharides, salts, enzymes, and antibodies (these are highly important, as the olfactory neurons provide a direct passage for infection to pass to the brain). This mucus acts as a solvent for odor molecules, flows constantly, and is replaced approximately every ten minutes.
In insects, smells are sensed by olfactory sensory neurons in the chemosensory sensilla, which are present in insect antenna, palps, and tarsa, but also on other parts of the insect body. Odorants penetrate into the cuticle pores of chemosensory sensilla and get in contact with insect odorant-binding proteins (OBPs) or Chemosensory proteins (CSPs), before activating the sensory neurons.
Receptor neuron
The binding of the ligand (odor molecule or odorant) to the receptor leads to an action potential in the receptor neuron, via a second messenger pathway, depending on the organism. In mammals, the odorants stimulate adenylate cyclase to synthesize cAMP via a G protein called Golf. cAMP, which is the second messenger here, opens a cyclic nucleotide-gated ion channel (CNG), producing an influx of cations (largely Ca2+ with some Na+) into the cell, slightly depolarising it. The Ca2+ in turn opens a Ca2+-activated chloride channel, leading to efflux of Cl−, further depolarizing the cell and triggering an action potential. Ca2+ is then extruded through a sodium-calcium exchanger. A calcium-calmodulin complex also acts to inhibit the binding of cAMP to the cAMP-dependent channel, thus contributing to olfactory adaptation.
The main olfactory system of some mammals also contains small subpopulations of olfactory sensory neurons that detect and transduce odors somewhat differently. Olfactory sensory neurons that use trace amine-associated receptors (TAARs) to detect odors use the same second messenger signaling cascade as do the canonical olfactory sensory neurons. Other subpopulations, such as those that express the receptor guanylyl cyclase GC-D (Gucy2d) or the soluble guanylyl cyclase Gucy1b2, use a cGMP cascade to transduce their odorant ligands. These distinct subpopulations (olfactory subsystems) appear specialized for the detection of small groups of chemical stimuli.
This mechanism of transduction is somewhat unusual, in that cAMP works by directly binding to the ion channel rather than through activation of protein kinase A. It is similar to the transduction mechanism for photoreceptors, in which the second messenger cGMP works by directly binding to ion channels, suggesting that maybe one of these receptors was evolutionarily adapted into the other. There are also considerable similarities in the immediate processing of stimuli by lateral inhibition.
Averaged activity of the receptor neurons can be measured in several ways. In vertebrates, responses to an odor can be measured by an electro-olfactogram or through calcium imaging of receptor neuron terminals in the olfactory bulb. In insects, one can perform electroantennography or calcium imaging within the olfactory bulb.
Olfactory bulb projections
Olfactory sensory neurons project axons to the brain within the olfactory nerve, (cranial nerveI). These nerve fibers, lacking myelin sheaths, pass to the olfactory bulb of the brain through perforations in the cribriform plate, which in turn projects olfactory information to the olfactory cortex and other areas. The axons from the olfactory receptors converge in the outer layer of the olfactory bulb within small (≈50 micrometers in diameter) structures called glomeruli. Mitral cells, located in the inner layer of the olfactory bulb, form synapses with the axons of the sensory neurons within glomeruli and send the information about the odor to other parts of the olfactory system, where multiple signals may be processed to form a synthesized olfactory perception. A large degree of convergence occurs, with 25,000 axons synapsing on 25 or so mitral cells, and with each of these mitral cells projecting to multiple glomeruli. Mitral cells also project to periglomerular cells and granular cells that inhibit the mitral cells surrounding it (lateral inhibition). Granular cells also mediate inhibition and excitation of mitral cells through pathways from centrifugal fibers and the anterior olfactory nuclei. Neuromodulators like acetylcholine, serotonin and norepinephrine all send axons to the olfactory bulb and have been implicated in gain modulation, pattern separation, and memory functions, respectively.
The mitral cells leave the olfactory bulb in the lateral olfactory tract, which synapses on five major regions of the cerebrum: the anterior olfactory nucleus, the olfactory tubercle, the amygdala, the piriform cortex, and the entorhinal cortex. The anterior olfactory nucleus projects, via the anterior commissure, to the contralateral olfactory bulb, inhibiting it. The piriform cortex has two major divisions with anatomically distinct organizations and functions. The anterior piriform cortex (APC) appears to be better at determining the chemical structure of the odorant molecules, and the posterior piriform cortex (PPC) has a strong role in categorizing odors and assessing similarities between odors (e.g. minty, woody, and citrus are odors that can, despite being highly variant chemicals, be distinguished via the PPC in a concentration-independent manner). The piriform cortex projects to the medial dorsal nucleus of the thalamus, which then projects to the orbitofrontal cortex. The orbitofrontal cortex mediates conscious perception of the odor. The three-layered piriform cortex projects to a number of thalamic and hypothalamic nuclei, the hippocampus and amygdala and the orbitofrontal cortex, but its function is largely unknown. The entorhinal cortex projects to the amygdala and is involved in emotional and autonomic responses to odor. It also projects to the hippocampus and is involved in motivation and memory. Odor information is stored in long-term memory and has strong connections to emotional memory. This is possibly due to the olfactory system's close anatomical ties to the limbic system and hippocampus, areas of the brain that have long been known to be involved in emotion and place memory, respectively.
Since any one receptor is responsive to various odorants, and there is a great deal of convergence at the level of the olfactory bulb, it may seem strange that human beings are able to distinguish so many different odors. It seems that a highly complex form of processing must be occurring; however, as it can be shown that, while many neurons in the olfactory bulb (and even the pyriform cortex and amygdala) are responsive to many different odors, half the neurons in the orbitofrontal cortex are responsive to only one odor, and the rest to only a few. It has been shown through microelectrode studies that each individual odor gives a particular spatial map of excitation in the olfactory bulb. It is possible that the brain is able to distinguish specific odors through spatial encoding, but temporal coding must also be taken into account. Over time, the spatial maps change, even for one particular odor, and the brain must be able to process these details as well.
Inputs from the two nostrils have separate inputs to the brain, with the result that, when each nostril takes up a different odorant, a person may experience perceptual rivalry in the olfactory sense akin to that of binocular rivalry.
In insects, smells are sensed by sensilla located on the antenna and maxillary palp and first processed by the antennal lobe (analogous to the olfactory bulb), and next by the mushroom bodies and lateral horn.
Coding and perception
The process by which olfactory information is coded in the brain to allow for proper perception is still being researched, and is not completely understood. When an odorant is detected by receptors, they in a sense break the odorant down, and then the brain puts the odorant back together for identification and perception. The odorant binds to receptors that recognize only a specific functional group, or feature, of the odorant, which is why the chemical nature of the odorant is important.
After binding the odorant, the receptor is activated and will send a signal to the glomeruli in the olfactory bulb. Each glomerulus receives signals from multiple receptors that detect similar odorant features. Because several receptor types are activated due to the different chemical features of the odorant, several glomeruli are activated as well. The signals from the glomeruli are transformed to a pattern of oscillations of neural activities of the mitral cells, the output neurons from the olfactory bulb. Olfactory bulb sends this pattern to the olfactory cortex. Olfactory cortex is thought to have associative memories, so that it resonates to this bulbar pattern when the odor object is recognized. The cortex sends centrifugal feedback to the bulb. This feedback could suppress bulbar responses to the recognized odor objects, causing olfactory adaptation to background odors, so that the newly arrived foreground odor objects could be singled out for better recognition. During odor search, feedback could also be used to enhance odor detection. The distributed code allows the brain to detect specific odors in mixtures of many background odors.
It is a general idea that the layout of brain structures corresponds to physical features of stimuli (called topographic coding), and similar analogies have been made in smell with concepts such as a layout corresponding to chemical features (called chemotopy) or perceptual features. While chemotopy remains a highly controversial concept, evidence exists for perceptual information implemented in the spatial dimensions of olfactory networks.
Accessory olfactory system
Many animals, including most mammals and reptiles, but not humans, have two distinct and segregated olfactory systems: a main olfactory system, which detects volatile stimuli, and an accessory olfactory system, which detects fluid-phase stimuli. Behavioral evidence suggests that these fluid-phase stimuli often function as pheromones, although pheromones can also be detected by the main olfactory system. In the accessory olfactory system, stimuli are detected by the vomeronasal organ, located in the vomer, between the nose and the mouth. Snakes use it to smell prey, sticking their tongue out and touching it to the organ. Some mammals make a facial expression called flehmen to direct stimuli to this organ.
The sensory receptors of the accessory olfactory system are located in the vomeronasal organ. As in the main olfactory system, the axons of these sensory neurons project from the vomeronasal organ to the accessory olfactory bulb, which in the mouse is located on the dorsal-posterior portion of the main olfactory bulb. Unlike in the main olfactory system, the axons that leave the accessory olfactory bulb do not project to the brain's cortex but rather to targets in the amygdala and bed nucleus of the stria terminalis, and from there to the hypothalamus, where they may influence aggression and mating behavior.
In insects
Insect olfaction refers to the function of chemical receptors that enable insects to detect and identify volatile compounds for foraging, predator avoidance, finding mating partners (via pheromones) and locating oviposition habitats. Thus, it is the most important sensation for insects. Most important insect behaviors must be timed perfectly which is dependent on what they smell and when they smell it. For example, smell is essential for hunting in many species of wasps, including Polybia sericea.
The two organs insects primarily use for detecting odors are the antennae and specialized mouth parts called the maxillary palps. However, a recent study has demonstrated the olfactory role of ovipositor in fig wasps. Inside of these olfactory organs there are neurons called olfactory receptor neurons which, as the name implies, house receptors for scent molecules in their cell membranes. The majority of olfactory receptor neurons typically reside in the antenna. These neurons can be very abundant, for example Drosophila flies have 2,600 olfactory sensory neurons.
Insects are capable of smelling and differentiating between thousands of volatile compounds both sensitively and selectively. Sensitivity is how attuned the insect is to very small amounts of an odorant or small changes in the concentration of an odorant. Selectivity refers to the insects' ability to tell one odorant apart from another. These compounds are commonly broken into three classes: short chain carboxylic acids, aldehydes and low molecular weight nitrogenous compounds. Some insects, such as the moth Deilephila elpenor, use smell as a means to find food sources.
In plants
The tendrils of plants are especially sensitive to airborne volatile organic compounds. Parasites such as dodder make use of this in locating their preferred hosts and locking on to them. The emission of volatile compounds is detected when foliage is browsed by animals. Threatened plants are then able to take defensive chemical measures, such as moving tannin compounds to their foliage.
Machine-based smelling
Scientists have devised methods for quantifying the intensity of odors, in particular for the purpose of analyzing unpleasant or objectionable odors released by an industrial source into a community. Since the 1800s industrial countries have encountered incidents where proximity of an industrial source or landfill produced adverse reactions among nearby residents regarding airborne odor. The basic theory of odor analysis is to measure what extent of dilution with "pure" air is required before the sample in question is rendered indistinguishable from the "pure" or reference standard. Since each person perceives odor differently, an "odor panel" composed of several different people is assembled, each sniffing the same sample of diluted specimen air. A field olfactometer can be utilized to determine the magnitude of an odor.
Many air management districts in the US have numerical standards of acceptability for the intensity of odor that is allowed to cross into a residential property. For example, the Bay Area Air Quality Management District has applied its standard in regulating numerous industries, landfills, and sewage treatment plants. Example applications this district has engaged are the San Mateo, California, wastewater treatment plant; the Shoreline Amphitheatre in Mountain View, California; and the IT Corporation waste ponds, Martinez, California.
Classification
Systems of classifying odors include:
Crocker-Henderson system, which rates smells on a 0-8 scale for each of four "primary" smells: fragrant, acid, burnt, and caprylic.
Henning's prism
Zwaardemaker smell system (invented by Hendrik Zwaardemaker)
Disorders
Specific terms are used to describe disorders associated with smelling:
Anosmia: inability to smell
Hyperosmia: an abnormally acute sense of smell
Hyposmia: decreased ability to smell
Presbyosmia: the natural decline in the sense of smell in old age
Dysosmia: distortion in the sense of smell
Parosmia: distortion in the perception of an odor
Phantosmia: distortion in the absence of an odor, "hallucinated smell"
Heterosmia: inability to distinguish odors
Desiderosmia: a compulsive craving for certain odors (e.g., rubber tires, gasoline). This has been associated with iron deficiency anemia.
Olfactory reference syndrome: psychological disorder that causes the patient to imagine he or she has strong body odor
Osmophobia: aversion or psychological hypersensitivity to odors
Primary Sjögren's syndrome: may impair chemosensory function, including taste and smell
Viruses can also infect the olfactory epithelium leading to a loss of the sense of olfaction. About 50% of patients with SARS-CoV-2 (causing COVID-19) experience some type of disorder associated with their sense of smell, including anosmia and parosmia. SARS-CoV-1, MERS-CoV and even the flu (influenza virus) can also disrupt olfaction.
| Biology and health sciences | Nervous system | null |
21244409 | https://en.wikipedia.org/wiki/VY%20Canis%20Majoris | VY Canis Majoris | VY Canis Majoris (abbreviated to VY CMa) is an extreme oxygen-rich red hypergiant or red supergiant (O-rich RHG or RSG) and pulsating variable star from the Solar System in the slightly southern constellation of Canis Major. It is one of the largest known stars, one of the most luminous and massive red supergiants, and one of the most luminous stars in the Milky Way.
No evidence has been found that it is part of a multiple star system. Its great infrared (IR) excess makes it one of the brightest objects in the local part of the galaxy at wavelengths of 5 to 20 microns (μm) and indicates a dust shell or heated disk. It is about times the mass of the Sun (). It is surrounded by a complex asymmetric circumstellar envelope (CSE) caused by its mass loss. It produces strong molecular maser emission and was one of the first radio masers discovered. VY CMa is embedded in the large molecular cloud Sh2-310, a large, quite local star-forming H II region—its diameter: 480 arcminutes (′) or .
The radius of VY CMa is 1,420 times that of the Sun (), which is close to the modelled maximum, the Hayashi limit, a volume nearly 3 billion times that of the Sun. Taking this mid-point estimate as correct, an object travelling at the speed of light would take 6 hours to go around its surface, compared to 14.5 seconds for the Sun. If this star replaced the Sun its surface would, per this approximation, be beyond the orbit of Jupiter.
Observational history
The first known-recorded observation of VY Canis Majoris is in the star catalogue of the French astronomer Jérôme Lalande in 1801, which lists it as a 7th order of magnitude star. Further quite frequent studies of its apparent magnitude imply the light of the star as viewed from Earth has faded since 1850, which could be due to emission changes or a denser part of its surrounds becoming interposed (extinction).
Since 1847, VY Canis Majoris has been described as a crimson star. During the 19th century, observers measured at least six discrete components, suggesting that it might be a multiple star. These are now known to be bright zones in the host nebula. Observations in 1957 and high-resolution imaging in 1998 all but rule out any companion stars.
Giving spectral lines in brackets, the star is a strong emitter of OH (1612 MHz), (22235.08 MHz), and (43122 MHz) masers, which has been proven to be typical of an OH/IR star. Molecules, such as , , , , , , , and have been detected.
The variation in the star's brightness was first described in 1931, when it was listed (in German) as a long-period variable with a photographic magnitude range of 9.5 to 11.5. It was given the variable star designation VY Canis Majoris in 1939, the 43rd variable star of the constellation Canis Major.
Combining data from the mentioned telescope with others from the Keck in Hawaii it was possible to make a three-dimensional reconstruction of the envelope of the star. This reconstruction showed that the star's mass loss is much more complex than expected for any red supergiant or hypergiant. It became clear that the bows and nodules appeared at different times; the jets are randomly oriented, which prompts suspicion they derive from explosions of active parts of the photosphere. The spectroscopy proves the jets move away from the star at different speeds, confirming multiple events and directions as with coronal mass ejections. Multiple asymmetric mass loss events and the ejection of the outermost material are deduced to have occurred within the last 500 to 1,000 years, while that of a knot near the star would be less than 100 years. The mass loss is due to strong convection in the tenuous outer layers of the star, associated with magnetic fields. Ejections are analogous to—but much larger than—coronal ejections of the Sun.
Distance
In 1976, Lada and Reid published observations of the bright-rimmed molecular cloud Sh2-310, which is 15″ east of the star. At its edge bordered by the bright rim, an abrupt decrease in the CO emission and an increase in brightness of the emission were observed, indicating possible destruction of molecular material and enhanced heating at the cloud-rim interface, respectively. They assumed the distance of the cloud is approximately equal to that of the stars, which are members of the open cluster NGC 2362, that ionize the rim. NGC 2362 could be anywhere in the ranges of (kpc) or (ly) away as determined from its color-magnitude diagram. This star is projected onto the tip of the cloud rim, strongly suggesting its association. Furthermore, all the vectors of velocity of Sh2-310 are very close to those of the star. There is thus a near-certain physical association of the star with Sh2-310 and with NGC 2362 in all standard models. Sh2-310 besides containing VY Canis Majoris and NGC 2362 also is host to the dark nebulae, LDN 1660, LDN 1664, and LDN 1667. Sh2-310 is also host to the stars Tau Canis Majoris which is the brightest member of NGC 2362, UW Canis Majoris and HD 58011 which along with VY Canis Majoris are thought to be probable sources of ionization of gases in Sh2-310. Sh2-310 itself is located on the outer edge of the Orion Arm of the Milky Way. Melnik and others later prefer a range centred on 1.2 kiloparsecs (about 3,900 light-years).
Distances can be calculated by measuring the change in position against very distant background objects as the telescope orbits the Sun. However, this star has a small parallax due to its distance, and standard visual observations have a margin of error too large for a hypergiant star with an extended CSE to be useful, for example, the Hipparcos Catalogue of 1997 gives a purely notional parallax of (mas), in which the "central" figure equates to (). Parallax can be measured to high accuracy from the observation of masers using a long baseline interferometry. In 2008, such observations of masers using VERA interferometry from the National Astronomical Observatory of Japan gave a parallax of , corresponding to a distance of (about ). In 2012, observations of masers using very-long-baseline interferometry (VLBI) from Very Long Baseline Array (VLBA) independently derived a parallax of , corresponding to a distance of (about ). These imply the cloud (Sh2-310) is less remote than thought or that VY CMa is a foreground object. The Gaia mission provides highly constrained parallaxes to some objects, but the data release 2 value of for VY CMa is not meaningful.
Variability
VY Canis Majoris is a variable star that varies from an apparent visual magnitude of 9.6 at minimum brightness to a magnitude of 6.5 at maximum with an estimated pulsational period of 956 days. In the General Catalogue of Variable Stars (GCVS) it is classed a semiregular variable of sub-type SRc, indicating a cool supergiant, although it is classed as a type LC slow irregular variable star in the American Association of Variable Star Observers (AAVSO) Variable Star Index. Other periods of 1,600 and 2,200 days have been derived.
VY CMa is sometimes considered as the prototype for a class of heavily mass-losing OH/IR supergiants, distinct from the more common asymptotic giant branch OH/IR stars.
Spectrum
The spectrum of VY Canis Majoris is that of a high-luminosity M-class star. The hydrogen lines, however, have P Cygni profiles fit for luminous blue variables. The spectrum is dominated by TiO bands whose strengths suggest a classification of M5. The H-alpha (Hα) line is not visible yet and there are unusual emission lines of neutral elements such as sodium and calcium. The luminosity class as determined from different spectral features varies from bright giant (II) to bright supergiant (Ia), with a compromise being given: as M5eIbp. Old classifications were confused by the interpretation of surrounding nebulosity as companion stars.
The present spectral classification system is inadequate to this star's complexities. The class depends on which of its complex spectral features are stressed. Further, key facets vary over time as to this star. It is cooler and thus redder than M2, and is usually classified between M3 and M5. A class as extreme as M2.5 appeared in a study of 2006. The luminosity class is likewise confused and often given only as I, partly because luminosity classes are poorly defined in the red and infrared portions of the spectrum. One study though, gives a luminosity class of Ia+ which means a hypergiant or extremely luminous supergiant.
Physical properties
A very large and luminous star, VY CMa is among the most extreme stars in the Milky Way and has an effective temperature below . It occupies the upper-right hand corner of the HR diagram although its exact luminosity and temperature are uncertain. Most of the properties of the star depend directly on its distance.
Luminosity
The bolometric luminosity (Lbol) of VY CMa can be calculated from spectral energy distribution or bolometric flux, which can be determined from photometry in several visible and infrared bands. Earlier calculations of the luminosity based on an assumed distance of gave luminosities between 200,000 and 560,000 times the Sun's luminosity (). This is considerably very close or beyond the empirical Humphreys–Davidson limit. One study gave nearly at a distance of . In 2006 a luminosity of was calculated by integrating the total fluxes over the entire nebula, since most of the radiation coming from the star is reprocessed by the dust in the surrounding cloud. More recent estimates of the luminosity extrapolate values below based on distances below 1.2 kpc.
Most of the output of VY CMa is emitted as infrared radiation, with a maximum emission at , which is in part caused by reprocessing of the radiation by the circumstellar nebula. Many older luminosity estimates are consistent with current ones if they are re-scaled to the distance of 1.2 kpc. Despite being one of the most luminous stars in the Milky Way, much of the visible light of VY CMa is absorbed by the circumstellar envelope, so the star needs a telescope to be observed. Removing its envelope, the star would be one for the naked eye.
Mass
Since this star has no companion star, its mass cannot be measured directly through gravitational interactions. Comparison of the effective temperature and bolometric luminosity compared to evolutionary tracks for massive stars suggest its initial mass was for a rotating star but current mass —or at first if non-rotating falling to present-day , and an age of 8.2 million years (Myr). Older studies have found much higher initial masses (thus also higher current masses) or a progenitor mass of based on old luminosity estimates.
Mass loss
VY CMa has a strong stellar wind and is losing much material due to its high luminosity and quite low surface gravity. It has an average mass loss rate of per year, among the highest known and unusually high even for a red supergiant, as evidenced by its extensive envelope. It is thus an exponent for the understanding of high-mass loss episodes near the end of massive star evolution. The mass loss rate probably exceeded /yr during the greatest mass loss events.
The star has produced large, probably convection-driven, mass-loss events 70, 120, 200, and 250 years ago. The clump shed by the star between 1985 and 1995 is the source of its hydroxyl maser emission.
Temperature
The effective temperature of this star is uncertain. Some signature changes in its spectrum correspond to temperature variations. Early estimates of the mean temperature assumed values below 3,000 K based on a spectral class of M5. In 2006, its temperature was calculated to be as high as , corresponding to a spectral class of M2.5, yet this star is usually considered as an M4 to M5 star. Adopting the latter classes with the temperature scale proposed by Emily Levesque gives a range of between 3,450 and 3,535 K.
Size
The calculation of the radius of VY CMa is complicated by the extensive circumstellar envelope of the star. VY CMa is also a pulsating star, so its size changes with time. Earlier direct measurements of the radius at infrared (K-band = 2.2 μm) wavelength gave an angular diameter of , corresponding to radii above at an assumed distance of 1.5 kpc, considerably larger than expected for any red supergiant or red hypergiant. However, this is probably larger than the actual size of the underlying star and the angular diameter estimate appears exceedingly large due to interference by the circumstellar envelope. In 2006–2007 radii of have been derived from the estimated luminosity of and temperatures of 3,450–3,535 K.
On 6 and 7 March 2011, VY CMa was observed at near-infrared wavelengths using interferometry at the Very Large Telescope. The size of the star was calculated using the Rosseland Radius, the location at which the optical depth is , with two modern distances of and . Its angular diameter was directly measured at , which corresponds to a radius of at a distance of . The high spectral resolution of these observations allowed the effects of contamination by circumstellar layers to be minimised. An effective temperature of , corresponding to a spectral class of M4, was then derived from the radius and a luminosity of which is based on the distance and a measured flux of . In late 2013, a radius of was determined, based on a rather cool adopted temperature of 2,800 K and a luminosity of .
Most radius estimates of the VY CMa are considered as the size for the optical photosphere while the size of the star for the radio photosphere is calculated to be twice that of the size of the star for the optical photosphere. Despite the mass and very large size (though some estimates give smaller sizes), VY CMa has an average density of 5.33 to 8.38 mg/m3 (0.00000533 to 0.00000838 kg/m3), it is over 100,000 times less dense than Earth's atmosphere at sea level (1.2 kg/m3).
Largest star
VY Canis Majoris has been known to be an extreme object since the middle of the 20th century, although its true nature was uncertain. In the late 20th century, it was accepted that it was a post-main sequence red supergiant. Its angular diameter had been measured and found to be significantly different depending on the observed wavelength. The first meaningful estimates of its properties showed a very large star.
Early direct measurements of the radius at infrared (K-band = 2.2 μm) wavelength gave an angular diameter of , corresponding to radii above at a still very plausible distance of 1.5 kiloparsecs; a radius dwarfing other known red hypergiants. However, this is probably larger than the actual size of the underlying star—this angular diameter estimate is heightened from interference by the envelope. In 2006–07, radius between has been derived from the preferring luminosity of and the still-preferred temperature range of .
In contrast to prevailing opinion, a 2006 study, ignoring the effects of the circumstellar envelope in the observed flux of the star, derived a luminosity of , suggesting an initial mass of and radius of based on an assumed effective temperature of 3,650 K and distance of . On this basis they considered VY CMa and another notable extreme cool hypergiant star, NML Cygni, as normal early-type red supergiants. They assert that earlier very high luminosities of and very large radii of (or even ) were based on effective temperatures below 3,000 K that were unreasonably low.
Almost immediately another paper published a size estimate of and concluded that VY CMa is a true hypergiant. This uses the later well-reviewed effective temperature , and a luminosity of based on SED integration and a distance of .
In 2011, the star was studied at near-infrared wavelengths using interferometry at the Very Large Telescope. The size of the star was published at its Rosseland Radius, outside of which optical depth falls below , given the mean of two most modern, similar but distinct distances. Its angular diameter was directly measured at , thus radius of given a distance of . The high spectral resolution of these observations allowed the effects of contamination by circumstellar layers to be minimised. An effective temperature of , corresponding to a spectral class of M4, was then derived from the radius and a luminosity of which is based on the distance and a measured flux of .
Most such radius estimates are considered as the size for the mean limit of the optical photosphere while the size of the star for the radio photosphere is calculated to be twice that. Despite the mass and very large size (though some estimates give smaller sizes), VY CMa has an average density of 5.33 to 8.38 mg/m3 (0.00000533 to 0.00000838 kg/m3). It is over 100,000 times less dense than Earth's atmosphere at sea level (1.2 kg/m3).
In 2012, the size was calculated more accurately to be somewhat lower, for example , which leaves larger sizes published and in-date for other galactic and extragalactic red supergiants (and hypergiants) such as Westerlund 1 W26 and WOH G64. Despite this, VY Canis Majoris is still often described as the largest known star, sometimes with caveats to account for the highly uncertain sizes of all these stars. A 2013 estimate based on the Wittkowski radius and the Monnier radius put mean size at , and later that year, Matsuura and others put forward a competing method of finding radius within the envelope, putting the star at , based on a cool-end of estimates adopted temperature of 2,800 K and a luminosity of . However, these values are not consistent with its spectral types, leaving the 2012 values in better match.
Surroundings
VY Canis Majoris is surrounded by an extensive and dense asymmetric red reflection nebula, with a total ejected mass of and a temperature of , based on a DUSTY model atmosphere that has been formed by material expelled from its central star. The inner shell figures as 0.12 ″ across, corresponding to for a star 1,200 parsecs away, whereas that of the outer one is at 10″, corresponding to . This nebula is so bright that it was discovered in a dry night sky in 1917 with an 18 cm telescope, and its condensations were once regarded as companion stars. It has been extensively studied with the aid of the Hubble Space Telescope (HST), showing that the nebula has a complex structure that includes filaments and arcs, which were caused by past eruptions; the structure is akin to that around the post-red supergiant yellow hypergiant (Post-RSG YHG) IRC +10420. The similarity has led at least two professional articles to propose a model that the star might evolve blueward on the Hertzsprung–Russell diagram (HR diagram) to become a yellow hypergiant, then a luminous blue variable (LBV), and finally a Wolf–Rayet star (WR star).
Evolution
VY Canis Majoris is a highly evolved star yet less than 10 million years old (Myr old). Some old writings envisaged the star as a very young protostar or a massive pre-main-sequence star with an age of only 1 Myr and typically a circumstellar disk. It has probably evolved from a hot, dense O9 main sequence star of (solar radii). The star has evolved rapidly because of its high mass. The time spent to the red hypergiant phase is estimated to be between 100,000 and 500,000 years, and thus VY CMa most likely left its main sequence phase more than a million years ago.
The future evolution of VY CMa is uncertain, but like the most cool supergiants, the star will certainly explode as a supernova. It has begun to fuse helium into carbon en masse. Like Betelgeuse, it is losing mass and is expected to explode as a supernova within the next 100,000 years — it will probably revert to a higher temperature beforehand. The star is very unstable, having a prodigious mass loss such as in ejections.
VY Canis Majoris is a candidate for a star in a second red supergiant phase, but this is mostly speculative and unconfirmed.
From this star CO emission is coincident with the bright KI shell in its asymmetric nebula.
The star will produce either:
a moderately luminous and long-lasting type IIn supernova (SN IIn)
a hypernova; or a
superluminous supernova (SLSN) comparable to SN 1988Z
or less likely, a type Ib supernova, but it is unlikely that would be as luminous as SN 2006tf or SN 2006gy.
The explosion could be associated with gamma-ray bursts (GRB), and it will produce a shock wave of a speed of a few thousand kilometers per second that could hit the surrounding envelope of material, causing strong emission for many years after the explosion. For a star so large, the remnant would be probably a black hole rather than a neutron star.
| Physical sciences | Notable stars | Astronomy |
21244416 | https://en.wikipedia.org/wiki/Hypergiant | Hypergiant | A hypergiant (luminosity class 0 or Ia+) is a very rare type of star that has an extremely high luminosity, mass, size and mass loss because of its extreme stellar winds. The term hypergiant is defined as luminosity class 0 (zero) in the MKK system. However, this is rarely seen in literature or in published spectral classifications, except for specific well-defined groups such as the yellow hypergiants, RSG (red supergiants), or blue B(e) supergiants with emission spectra. More commonly, hypergiants are classed as Ia-0 or Ia+, but red supergiants are rarely assigned these spectral classifications. Astronomers are interested in these stars because they relate to understanding stellar evolution, especially star formation, stability, and their expected demise as supernovae. Notable examples of hypergiants include the Pistol Star, a blue hypergiant located close to the Galactic Center and one of the most luminous stars known; Rho Cassiopeiae, a yellow hypergiant that is one of the brightest to the naked eye; and Mu Cephei (Herschel's "Garnet Star"), one of the largest and brightest stars known.
Origin and definition
In 1956, the astronomers Feast and Thackeray used the term super-supergiant (later changed into hypergiant) for stars with an absolute magnitude brighter than MV = −7 (MBol will be larger for very cool and very hot stars, for example at least −9.7 for a B0 hypergiant). In 1971, Keenan suggested that the term would be used only for supergiants showing at least one broad emission component in Hα, indicating an extended stellar atmosphere or a relatively large mass loss rate. The Keenan criterion is the one most commonly used by scientists today; hence it is possible for a supergiant star to have a higher luminosity than a hypergiant of the same spectral class.
Hypergiants are expected to have a characteristic broadening and red-shifting of their spectral lines, producing a distinctive spectral shape known as a P Cygni profile. The use of hydrogen emission lines is not helpful for defining the coolest hypergiants, and these are largely classified by luminosity since mass loss is almost inevitable for the class.
Formation
Stars with an initial mass above about quickly move away from the main sequence and increase somewhat in luminosity to become blue supergiants. They cool and enlarge at approximately constant luminosity to become a red supergiant, then contract and increase in temperature as the outer layers are blown away. They may "bounce" backwards and forwards executing one or more "blue loops", still at a fairly steady luminosity, until they explode as a supernova or completely shed their outer layers to become a Wolf–Rayet star. Stars with an initial mass above about are simply too luminous to develop a stable extended atmosphere and so they never cool sufficiently to become red supergiants. The most massive stars, especially rapidly rotating stars with enhanced convection and mixing, may skip these steps and move directly to the Wolf–Rayet stage.
This means that stars at the top of the Hertzsprung–Russell diagram where hypergiants are found may be newly evolved from the main sequence and still with high mass, or much more evolved post-red supergiant stars that have lost a significant fraction of their initial mass, and these objects cannot be distinguished simply on the basis of their luminosity and temperature. High-mass stars with a high proportion of remaining hydrogen are more stable, while older stars with lower masses and a higher proportion of heavy elements have less stable atmospheres due to increased radiation pressure and decreased gravitational attraction. These are thought to be the hypergiants, near the Eddington limit and rapidly losing mass.
The yellow hypergiants are thought to be generally post-red supergiant stars that have already lost most of their atmospheres and hydrogen. A few more stable high mass yellow supergiants with approximately the same luminosity are known and thought to be evolving towards the red supergiant phase, but these are rare as this is expected to be a rapid transition. Because yellow hypergiants are post-red supergiant stars, there is a fairly hard upper limit to their luminosity at around , but blue hypergiants can be much more luminous, sometimes several million .
Almost all hypergiants exhibit variations in luminosity over time due to instabilities within their interiors, but these are small except for two distinct instability regions where luminous blue variables (LBVs) and yellow hypergiants are found. Because of their high masses, the lifetime of a hypergiant is very short in astronomical timescales: only a few million years compared to around 10 billion years for stars like the Sun. Hypergiants are only created in the largest and densest areas of star formation and because of their short lives, only a small number are known despite their extreme luminosity that allows them to be identified even in neighbouring galaxies. The time spent in some phases such as LBVs can be as short as a few thousand years.
Stability
As the luminosity of stars increases greatly with mass, the luminosity of hypergiants often lies very close to the Eddington limit, which is the luminosity at which the radiation pressure expanding the star outward equals the force of the star's gravity collapsing the star inward. This means that the radiative flux passing through the photosphere of a hypergiant may be nearly strong enough to lift off the photosphere. Above the Eddington limit, the star would generate so much radiation that parts of its outer layers would be thrown off in massive outbursts; this would effectively restrict the star from shining at higher luminosities for longer periods.
A good candidate for hosting a continuum-driven wind is Eta Carinae, one of the most massive stars ever observed. With an estimated mass of around 130 solar masses and a luminosity four million times that of the Sun, astrophysicists speculate that Eta Carinae may occasionally exceed the Eddington limit. The last time might have been a series of outbursts observed in 1840–1860, reaching mass loss rates much higher than our current understanding of what stellar winds would allow.
As opposed to line-driven stellar winds (that is, ones driven by absorbing light from the star in huge numbers of narrow spectral lines), continuum driving does not require the presence of "metallic" atoms — atoms other than hydrogen and helium, which have few such lines — in the photosphere. This is important, since most massive stars also are very metal-poor, which means that the effect must work independently of the metallicity. In the same line of reasoning, the continuum driving may also contribute to an upper mass limit even for the first generation of stars right after the Big Bang, which did not contain any metals at all.
Another theory to explain the massive outbursts of, for example, Eta Carinae is the idea of a deeply situated hydrodynamic explosion, blasting off parts of the star's outer layers. The idea is that the star, even at luminosities below the Eddington limit, would have insufficient heat convection in the inner layers, resulting in a density inversion potentially leading to a massive explosion. The theory has, however, not been explored very much, and it is uncertain whether this really can happen.
Another theory associated with hypergiant stars is the potential to form a pseudo-photosphere, that is a spherical optically dense surface that is actually formed by the stellar wind rather than being the true surface of the star. Such a pseudo-photosphere would be significantly cooler than the deeper surface below the outward-moving dense wind. This has been hypothesized to account for the "missing" intermediate-luminosity LBVs and the presence of yellow hypergiants at approximately the same luminosity and cooler temperatures. The yellow hypergiants are actually the LBVs having formed a pseudo-photosphere and so apparently having a lower temperature.
Relationships with Ofpe, WNL, LBV, and other supergiant stars
Hypergiants are evolved, high luminosity, high-mass stars that occur in the same or similar regions of the Hertzsprung–Russell diagram as some stars with different classifications. It is not always clear whether the different classifications represent stars with different initial conditions, stars at different stages of an evolutionary track, or is just an artifact of our observations. Astrophysical models explaining the phenomena show many areas of agreement. Yet there are some distinctions that are not necessarily helpful in establishing relationships between different types of stars.
Although most supergiant stars are less luminous than hypergiants of similar temperature, a few fall within the same luminosity range. Ordinary supergiants compared to hypergiants often lack the strong hydrogen emissions whose broadened spectral lines indicate significant mass loss. Evolved lower mass supergiants do not return from the red supergiant phase, either exploding as supernovae or leaving behind a white dwarf.
Luminous blue variables are a class of highly luminous hot stars that display characteristic spectral variation. They often lie in a "quiescent" zone with hotter stars generally being more luminous, but periodically undergo large surface eruptions and move to a narrow zone where stars of all luminosities have approximately the same temperature, around . This "active" zone is near the hot edge of the unstable "void" where yellow hypergiants are found, with some overlap. It is not clear whether yellow hypergiants ever manage to get past the instability void to become LBVs or explode as a supernova.
Blue hypergiants are found in the same parts of the HR diagram as LBVs but do not necessarily show the LBV variations. Some but not all LBVs show the characteristics of hypergiant spectra at least some of the time, but many authors would exclude all LBVs from the hypergiant class and treat them separately. Blue hypergiants that do not show LBV characteristics may be progenitors of LBVs, or vice versa, or both. Lower mass LBVs may be a transitional stage to or from cool hypergiants or are different type of object.
Wolf–Rayet stars are extremely hot stars that have lost much or all of their outer layers. WNL is a term used for late stage (i.e. cooler) Wolf–Rayet stars with spectra dominated by nitrogen. Although these are generally thought to be the stage reached by hypergiant stars after sufficient mass loss, it is possible that a small group of hydrogen-rich WNL stars are actually progenitors of blue hypergiants or LBVs. These are the closely related Ofpe (O-type spectra plus H, He, and N emission lines, and other peculiarities) and WN9 (the coolest nitrogen Wolf–Rayet stars) which may be a brief intermediate stage between high mass main-sequence stars and hypergiants or LBVs. Quiescent LBVs have been observed with WNL spectra and apparent Ofpe/WNL stars have changed to show blue hypergiant spectra. High rotation rates cause massive stars to shed their atmospheres quickly and prevent the passage from main sequence to supergiant, so these directly become Wolf–Rayet stars. Wolf Rayet stars, slash stars, cool slash stars (aka WN10/11), Ofpe, Of+, and Of* stars are not considered hypergiants. Although they are luminous and often have strong emission lines, they have characteristic spectra of their own.
Known hypergiants
Hypergiants are difficult to study due to their rarity. Many hypergiants have highly variable spectra, but they are grouped here into broad spectral classes.
Luminous blue variables
Some luminous blue variables are classified as hypergiants, during at least part of their cycle of variation:
AG Carinae, a massive Luminous blue variable and a part of the Carina constellation, which is transitioning from an O-type star to a Wolf-Rayet star.
Eta Carinae, inside the Carina Nebula (NGC 3372) in the southern constellation of Carina. Eta Carinae is extremely massive, possibly as much as 120 to 150 times the mass of the Sun, and is four to five million times as luminous. Possibly a different type of object from the LBVs, or extreme for a LBV.
P Cygni, in the northern constellation of Cygnus. Prototype for the general characteristics of LBV spectral lines.
S Doradus, in the Large Magellanic Cloud, in the southern constellation of Dorado. Prototype variable, LBVs are still sometimes called S Doradus variables.
The Pistol Star (V4647 Sgr), near the center of the Milky Way, in the constellation of Sagittarius. The Pistol Star is over 25 times more massive than the Sun, and is about 1.7 million times more luminous. Considered a candidate LBV, but variability has not been confirmed.
V4029 Sagittarii
V905 Scorpii
HD 6884, (R40 in SMC)
HD 269700, (R116 in the LMC)
LBV 1806-20 in the 1806-20 cluster on the other side of the Milky Way.
Blue hypergiants
Usually B-class, occasionally late O or early A:
2dFS 3235
AzV 2
AzV 65
AzV 76
AzV 78
AzV 367
Barbá 2-2
BP Crucis (Wray 977 or GX 301-2), binary with a pulsar companion.
Cygnus OB2-12
HD 5291 (Sk 56)
HD 32034 (R62 in LMC)
HD 37974 (R126 in LMC)
HD 80077, LBV candidate`
HD 268835 (R66 in LMC)
HD 269781 (in LMC)
HD 269661 (R111 in LMC)
HD 269604 (in LMC)
HDE 269128 (R81 in LMC), LBV candidate, eclipsing binary system.
HD 269896
HT Sagittae
M33 OB21 108
MAC 1-277
V430 Scuti
V452 Scuti, LBV candidate
V1429 Aquilae (= MWC 314), LBV candidate with a supergiant companion.
V1768 Cygni
V2140 Cygni
V4030 Sagittarii
6 Cassiopeiae
Zeta¹ Scorpii
In Galactic Center Region:
Star 13, type O, LBV candidate
Star 18, type O, LBV candidate
In Westerlund 1:
W5 (possible Wolf–Rayet)
W7
W13 (binary?)
W16a
W27
W30
W33
W42a
Yellow hypergiants
Yellow hypergiants typically have late A to early K spectra. However, A-type hypergiants can also be called white hypergiants.
HD 7583 (R45 in SMC)
HD 33579 (in LMC)
HD 268757 (R59 in LMC)
IRAS 17163-3907
IRAS 18357-0604
IRC+10420 (V1302 Aql)
Omicron1 Centauri
Rho Cassiopeiae
V382 Carinae
V509 Cassiopeiae
V766 Centauri (HR 5171A, possible red supergiant)
V810 Centauri A
V1427 Aquilae
V915 Scorpii
R Puppis
Variable A (in M33)
RW Cephei
In the Large Magellanic Cloud
WOH G64
In Westerlund 1:
W4
W8a
W12a
W16a
W32
W265
In the Sextans galaxy:
Sextans A7
In the LS1 galaxy/globular cluster:
Mothra (star)
Plus at least two probable cool hypergiants in the recently discovered Scutum Red Supergiant Clusters: F15 and possibly F13 in RSGC1 and Star 49 in RSGC2.
Red hypergiants
K to M type spectra, the largest known stars by radius. Hypergiant luminosity classes are rarely applied to red supergiants, although the term red hypergiant is sometimes applied to the most extended and unstable red supergiants, with radii on the order of .
Mu Cephei
VV Cephei A
NML Cygni
S Persei
VY Canis Majoris - potentially the largest star in the Milky Way
KY Cygni
PZ Cassiopeiae
HD 143183
UY Scuti
V602 Carinae
| Physical sciences | Stellar astronomy | Astronomy |
21245414 | https://en.wikipedia.org/wiki/Line%20integral | Line integral | In mathematics, a line integral is an integral where the function to be integrated is evaluated along a curve. The terms path integral, curve integral, and curvilinear integral are also used; contour integral is used as well, although that is typically reserved for line integrals in the complex plane.
The function to be integrated may be a scalar field or a vector field. The value of the line integral is the sum of values of the field at all points on the curve, weighted by some scalar function on the curve (commonly arc length or, for a vector field, the scalar product of the vector field with a differential vector in the curve). This weighting distinguishes the line integral from simpler integrals defined on intervals. Many simple formulae in physics, such as the definition of work as have natural continuous analogues in terms of line integrals, in this case which computes the work done on an object moving through an electric or gravitational field along a path
Vector calculus
In qualitative terms, a line integral in vector calculus can be thought of as a measure of the total effect of a given tensor field along a given curve. For example, the line integral over a scalar field (rank 0 tensor) can be interpreted as the area under the field carved out by a particular curve. This can be visualized as the surface created by and a curve C in the xy plane. The line integral of f would be the area of the "curtain" created—when the points of the surface that are directly over C are carved out.
Line integral of a scalar field
Definition
For some scalar field where , the line integral along a piecewise smooth curve is defined as
where is an arbitrary bijective parametrization of the curve such that and give the endpoints of and . Here, and in the rest of the article, the absolute value bars denote the standard (Euclidean) norm of a vector.
The function is called the integrand, the curve is the domain of integration, and the symbol may be intuitively interpreted as an elementary arc length of the curve (i.e., a differential length of ). Line integrals of scalar fields over a curve do not depend on the chosen parametrization of .
Geometrically, when the scalar field is defined over a plane , its graph is a surface in space, and the line integral gives the (signed) cross-sectional area bounded by the curve and the graph of . See the animation to the right.
Derivation
For a line integral over a scalar field, the integral can be constructed from a Riemann sum using the above definitions of , and a parametrization of . This can be done by partitioning the interval into sub-intervals of length , then denotes some point, call it a sample point, on the curve . We can use the set of sample points to approximate the curve as a polygonal path by introducing the straight line piece between each of the sample points and . (The approximation of a curve to a polygonal path is called rectification of a curve, see here for more details.) We then label the distance of the line segment between adjacent sample points on the curve as . The product of and can be associated with the signed area of a rectangle with a height and width of and , respectively. Taking the limit of the sum of the terms as the length of the partitions approaches zero gives us
By the mean value theorem, the distance between subsequent points on the curve, is
Substituting this in the above Riemann sum yields
which is the Riemann sum for the integral
Line integral of a vector field
Definition
For a vector field , the line integral along a piecewise smooth curve , in the direction of r, is defined as
where is the dot product, and is a regular parametrization (i.e: ) of the curve C such that and give the endpoints of C.
A line integral of a scalar field is thus a line integral of a vector field, where the vectors are always tangential to the line of the integration.
Line integrals of vector fields are independent of the parametrization r in absolute value, but they do depend on its orientation. Specifically, a reversal in the orientation of the parametrization changes the sign of the line integral.
From the viewpoint of differential geometry, the line integral of a vector field along a curve is the integral of the corresponding 1-form under the musical isomorphism (which takes the vector field to the corresponding covector field), over the curve considered as an immersed 1-manifold.
Derivation
The line integral of a vector field can be derived in a manner very similar to the case of a scalar field, but this time with the inclusion of a dot product. Again using the above definitions of , and its parametrization , we construct the integral from a Riemann sum. We partition the interval (which is the range of the values of the parameter ) into intervals of length . Letting be the th point on , then gives us the position of the th point on the curve. However, instead of calculating up the distances between subsequent points, we need to calculate their displacement vectors, . As before, evaluating at all the points on the curve and taking the dot product with each displacement vector gives us the infinitesimal contribution of each partition of on . Letting the size of the partitions go to zero gives us a sum
By the mean value theorem, we see that the displacement vector between adjacent points on the curve is
Substituting this in the above Riemann sum yields
which is the Riemann sum for the integral defined above.
Path independence
If a vector field is the gradient of a scalar field (i.e. if is conservative), that is,
then by the multivariable chain rule the derivative of the composition of and is
which happens to be the integrand for the line integral of on . It follows, given a path C, that
In other words, the integral of over C depends solely on the values of at the points and , and is thus independent of the path between them. For this reason, a line integral of a conservative vector field is called path independent.
Applications
The line integral has many uses in physics. For example, the work done on a particle traveling on a curve C inside a force field represented as a vector field is the line integral of on C.
Flow across a curve
For a vector field , , the line integral across a curve C ⊂ U, also called the flux integral, is defined in terms of a piecewise smooth parametrization , , as:
Here is the dot product, and is the clockwise perpendicular of the velocity vector
The flow is computed in an oriented sense: the curve has a specified forward direction from to , and the flow is counted as positive when is on the clockwise side of the forward velocity vector .
Complex line integral
In complex analysis, the line integral is defined in terms of multiplication and addition of complex numbers. Suppose U is an open subset of the complex plane C, is a function, and is a curve of finite length, parametrized by , where . The line integral
may be defined by subdividing the interval [a, b] into and considering the expression
The integral is then the limit of this Riemann sum as the lengths of the subdivision intervals approach zero.
If the parametrization is continuously differentiable, the line integral can be evaluated as an integral of a function of a real variable:
When is a closed curve (initial and final points coincide), the line integral is often denoted sometimes referred to in engineering as a cyclic integral.
To establish a complete analogy with the line integral of a vector field, one must go back to the definition of differentiability in multivariable calculus. The gradient is defined from Riesz representation theorem, and inner products in complex analysis involve conjugacy (the gradient of a function at some would be , and the complex inner product would attribute twice a conjugate to in the vector field definition of a line integral).
The line integral with respect to the conjugate complex differential is defined to be
The line integrals of complex functions can be evaluated using a number of techniques. The most direct is to split into real and imaginary parts, reducing the problem to evaluating two real-valued line integrals. The Cauchy integral theorem may be used to equate the line integral of an analytic function to the same integral over a more convenient curve. It also implies that over a closed curve enclosing a region where is analytic without singularities, the value of the integral is simply zero, or in case the region includes singularities, the residue theorem computes the integral in terms of the singularities. This also implies the path independence of complex line integral for analytic functions.
Example
Consider the function , and let the contour L be the counterclockwise unit circle about 0, parametrized by with in using the complex exponential. Substituting, we find:
This is a typical result of Cauchy's integral formula and the residue theorem.
Relation of complex line integral and line integral of vector field
Viewing complex numbers as 2-dimensional vectors, the line integral of a complex-valued function has real and complex parts equal to the line integral and the flux integral of the vector field corresponding to the conjugate function Specifically, if parametrizes L, and corresponds to the vector field then:
By Cauchy's theorem, the left-hand integral is zero when is analytic (satisfying the Cauchy–Riemann equations) for any smooth closed curve L. Correspondingly, by Green's theorem, the right-hand integrals are zero when is irrotational (curl-free) and incompressible (divergence-free). In fact, the Cauchy-Riemann equations for are identical to the vanishing of curl and divergence for .
By Green's theorem, the area of a region enclosed by a smooth, closed, positively oriented curve is given by the integral This fact is used, for example, in the proof of the area theorem.
Quantum mechanics
The path integral formulation of quantum mechanics actually refers not to path integrals in this sense but to functional integrals, that is, integrals over a space of paths, of a function of a possible path. However, path integrals in the sense of this article are important in quantum mechanics; for example, complex contour integration is often used in evaluating probability amplitudes in quantum scattering theory.
| Mathematics | Multivariable and vector calculus | null |
21245707 | https://en.wikipedia.org/wiki/Red%20giant | Red giant | A red giant is a luminous giant star of low or intermediate mass (roughly 0.3–8 solar masses ()) in a late phase of stellar evolution. The outer atmosphere is inflated and tenuous, making the radius large and the surface temperature around or lower. The appearance of the red giant is from yellow-white to reddish-orange, including the spectral types K and M, sometimes G, but also class S stars and most carbon stars.
Red giants vary in the way by which they generate energy:
most common red giants are stars on the red-giant branch (RGB) that are still fusing hydrogen into helium in a shell surrounding an inert helium core
red-clump stars in the cool half of the horizontal branch, fusing helium into carbon in their cores via the triple-alpha process
asymptotic-giant-branch (AGB) stars with a helium burning shell outside a degenerate carbon–oxygen core, and a hydrogen-burning shell just beyond that.
Many of the well-known bright stars are red giants because they are luminous and moderately common. The K0 RGB star Arcturus is 36 light-years away, and Gacrux is the nearest M-class giant at 88 light-years' distance.
A red giant will usually produce a planetary nebula and become a white dwarf at the end of its life.
Characteristics
A red giant is a star that has exhausted the supply of hydrogen in its core and has begun thermonuclear fusion of hydrogen in a shell surrounding the core. They have radii tens to hundreds of times larger than that of the Sun. However, their outer envelope is lower in temperature, giving them a yellowish-orange hue. Despite the lower energy density of their envelope, red giants are many times more luminous than the Sun because of their great size. Red-giant-branch stars have luminosities up to nearly three thousand times that of the Sun (); spectral types of K or M have surface temperatures of (compared with the Sun's photosphere temperature of nearly ) and radii up to about 200 times the Sun (). Stars on the horizontal branch are hotter, with only a small range of luminosities around . Asymptotic-giant-branch stars range from similar luminosities as the brighter stars of the red-giant branch, up to several times more luminous at the end of the thermal pulsing phase.
Among the asymptotic-giant-branch stars belong the carbon stars of type C-N and late C-R, produced when carbon and other elements are convected to the surface in what is called a dredge-up. The first dredge-up occurs during hydrogen shell burning on the red-giant branch, but does not produce a large carbon abundance at the surface. The second, and sometimes third, dredge-up occurs during helium shell burning on the asymptotic-giant branch and convects carbon to the surface in sufficiently massive stars.
The stellar limb of a red giant is not sharply defined, contrary to their depiction in many illustrations. Rather, due to the very low mass density of the envelope, such stars lack a well-defined photosphere, and the body of the star gradually transitions into a 'corona'. The coolest red giants have complex spectra, with molecular lines, emission features, and sometimes masers, particularly from thermally pulsing AGB stars. Observations have also provided evidence of a hot chromosphere above the photosphere of red giants, where investigating the heating mechanisms for the chromospheres to form requires 3D simulations of red giants.
Another noteworthy feature of red giants is that, unlike Sun-like stars whose photospheres have a large number of small convection cells (solar granules), red-giant photospheres, as well as those of red supergiants, have just a few large cells, the features of which cause the variations of brightness so common on both types of stars.
Evolution
Red giants are evolved from main-sequence stars with masses in the range from about to around . When a star initially forms from a collapsing molecular cloud in the interstellar medium, it contains primarily hydrogen and helium, with trace amounts of "metals" (in astrophysics, this refers to all elements heavier than hydrogen and helium). These elements are all uniformly mixed throughout the star. The star "enters" the main sequence when its core reaches a temperature (several million kelvins) high enough to begin fusing hydrogen-1 (the predominant isotope), and establishes hydrostatic equilibrium. (In astrophysics, stellar fusion is often referred to as "burning", with hydrogen fusion sometimes termed "hydrogen burning".) Over its main sequence life, the star slowly fuses the hydrogen in the core into helium; its main-sequence life ends when nearly all the hydrogen in the core has been fused. For the Sun, the main-sequence lifetime is approximately 10 billion years. More massive stars burn disproportionately faster and so have a shorter lifetime than less massive stars.
When the star has mostly exhausted the hydrogen fuel in its core, the core's rate of nuclear reactions declines, and thus so do the radiation and thermal pressure the core generates, which are what support the star against gravitational contraction. The star further contracts, increasing the pressures and thus temperatures inside the star (as described by the ideal gas law). Eventually a "shell" layer around the core reaches temperatures sufficient to fuse hydrogen and thus generate its own radiation and thermal pressure, which "re-inflates" the star's outer layers and causes them to expand. The hydrogen-burning shell results in a situation that has been described as the mirror principle: when the core within the shell contracts, the layers of the star outside the shell must expand. The detailed physical processes that cause this are complex. Still, the behavior is necessary to satisfy simultaneous conservation of gravitational and thermal energy in a star with the shell structure. The core contracts and heats up due to the lack of fusion, and so the outer layers of the star expand greatly, absorbing most of the extra energy from shell fusion. This process of cooling and expanding is the subgiant stage. When the envelope of the star cools sufficiently it becomes convective, the star stops expanding, its luminosity starts to increase, and the star is ascending the red-giant branch of the Hertzsprung–Russell (H–R) diagram.
The evolutionary path the star takes as it moves along the red-giant branch depends on the mass of the star. For the Sun and stars of less than about the core will become dense enough that electron degeneracy pressure will prevent it from collapsing further. Once the core is degenerate, it will continue to heat until it reaches a temperature of roughly , hot enough to begin fusing helium to carbon via the triple-alpha process. Once the degenerate core reaches this temperature, the entire core will begin helium fusion nearly simultaneously in a so-called helium flash. In more-massive stars, the collapsing core will reach these temperatures before it is dense enough to be degenerate, so helium fusion will begin much more smoothly, and produce no helium flash. The core helium fusing phase of a star's life is called the horizontal branch in metal-poor stars, so named because these stars lie on a nearly horizontal line in the H–R diagram of many star clusters. Metal-rich helium-fusing stars instead lie on the so-called red clump in the H–R diagram.
An analogous process occurs when the core helium is exhausted, and the star collapses once again, causing helium in a shell to begin fusing. At the same time, hydrogen may begin fusion in a shell just outside the burning helium shell. This puts the star onto the asymptotic giant branch, a second red-giant phase. The helium fusion results in the build-up of a carbon–oxygen core. A star below about will never start fusion in its degenerate carbon–oxygen core. Instead, at the end of the asymptotic-giant-branch phase the star will eject its outer layers, forming a planetary nebula with the core of the star exposed, ultimately becoming a white dwarf. The ejection of the outer mass and the creation of a planetary nebula finally ends the red-giant phase of the star's evolution. The red-giant phase typically lasts only around a billion years in total for a solar mass star, almost all of which is spent on the red-giant branch. The horizontal-branch and asymptotic-giant-branch phases proceed tens of times faster.
If the star has about 0.2 to , it is massive enough to become a red giant but does not have enough mass to initiate the fusion of helium. These "intermediate" stars cool somewhat and increase their luminosity but never achieve the tip of the red-giant branch and helium core flash. When the ascent of the red-giant branch ends they puff off their outer layers much like a post-asymptotic-giant-branch star and then become a white dwarf.
Stars that do not become red giants
Very-low-mass stars are fully convective and may continue to fuse hydrogen into helium for up to a trillion years until only a small fraction of the entire star is hydrogen. Luminosity and temperature steadily increase during this time, just as for more-massive main-sequence stars, but the length of time involved means that the temperature eventually increases by about 50% and the luminosity by around 10 times. Eventually the level of helium increases to the point where the star ceases to be fully convective and the remaining hydrogen locked in the core is consumed in only a few billion more years. Depending on mass, the temperature and luminosity continue to increase for a time during hydrogen shell burning, the star can become hotter than the Sun and tens of times more luminous than when it formed although still not as luminous as the Sun. After some billions more years, they start to become less luminous and cooler even though hydrogen shell burning continues. These become cool helium white dwarfs.
Very-high-mass stars develop into supergiants that follow an evolutionary track that takes them back and forth horizontally over the H–R diagram, at the right end constituting red supergiants. These usually end their life as a type II supernova. The most massive stars can become Wolf–Rayet stars without becoming giants or supergiants at all.
Planets
Prospects for habitability
Although traditionally it has been suggested the evolution of a star into a red giant will render its planetary system, if present, uninhabitable, some research suggests that, during the evolution of a star along the red-giant branch, it could harbor a habitable zone for several billion years at 2 astronomical units (AU) out to around 100 million years at out, giving perhaps enough time for life to develop on a suitable world. After the red-giant stage, there would for such a star be a habitable zone between for an additional one billion years. Later studies have refined this scenario, showing how for a star the habitable zone lasts from 100 million years for a planet with an orbit similar to that of Mars to 210 million years for one that orbits at Saturn distance to the Sun, the maximum time (370 million years) corresponding for planets orbiting at the distance of Jupiter. However, planets orbiting a star in equivalent orbits to those of Jupiter and Saturn would be in the habitable zone for 5.8 billion years and 2.1 billion years, respectively; for stars more massive than the Sun, the times are considerably shorter.
Enlargement of planets
As of 2023, several hundred giant planets have been discovered around giant stars. However, these giant planets are more massive than the giant planets found around solar-type stars. This could be because giant stars are more massive than the Sun (less massive stars will still be on the main sequence and will not have become giants yet) and more massive stars are expected to have more massive planets. However, the masses of the planets that have been found around giant stars do not correlate with the masses of the stars; therefore, the planets could be growing in mass during the stars' red giant phase. The growth in planet mass could be partly due to accretion from stellar wind, although a much larger effect would be Roche lobe overflow causing mass-transfer from the star to the planet when the giant expands out to the orbital distance of the planet. (A similar process in multiple star systems is believed to be the cause of most novas and type Ia supernovas.)
Examples
Many of the well-known bright stars are red giants, because they are luminous and moderately common. The red-giant branch variable star Gamma Crucis is the nearest M-class giant star at 88 light-years. The K1.5 red-giant branch star Arcturus is 36 light-years away.
Red-giant branch
Aldebaran (α Tauri)
Arcturus (α Bootis)
μ Leonis
Gacrux (γ Crucis)
Red-clump giants
Pollux (β Geminorum)
Capella Aa (α Aurigae)
α Cassiopeiae (Schedar)
δ Andromedae
Asymptotic giant branch
ρ Persei (Gorgonea Tertia)
Mira (ο Ceti)
χ Cygni
α Herculis (Rasalgethi)
The Sun as a red giant
The Sun will exit the main sequence in approximately 5 billion years and start to turn into a red giant. As a red giant, the Sun will grow so large (over 200 times its present-day radius: ; ) that it will engulf Mercury, Venus, and likely Earth. It will lose 38% of its mass growing, then will die into a white dwarf.
| Physical sciences | Stellar astronomy | null |
886766 | https://en.wikipedia.org/wiki/Bell%20test | Bell test | A Bell test, also known as Bell inequality test or Bell experiment, is a real-world physics experiment designed to test the theory of quantum mechanics in relation to Albert Einstein's concept of local realism. Named for John Stewart Bell, the experiments test whether or not the real world satisfies local realism, which requires the presence of some additional local variables (called "hidden" because they are not a feature of quantum theory) to explain the behavior of particles like photons and electrons. The test empirically evaluates the implications of Bell's theorem. , all Bell tests have found that the hypothesis of local hidden variables is inconsistent with the way that physical systems behave.
Many types of Bell tests have been performed in physics laboratories, often with the goal of ameliorating problems of experimental design or set-up that could in principle affect the validity of the findings of earlier Bell tests. This is known as "closing loopholes in Bell tests".
Bell inequality violations are also used in some quantum cryptography protocols, whereby a spy's presence is detected when Bell's inequalities cease to be violated.
Overview
The Bell test has its origins in the debate between Einstein and other pioneers of quantum physics, principally Niels Bohr. One feature of the theory of quantum mechanics under debate was the meaning of Heisenberg's uncertainty principle. This principle states that if some information is known about a given particle, there is some other information about it that is impossible to know. An example of this is found in observations of the position and the momentum of a given particle. According to the uncertainty principle, a particle's momentum and its position cannot simultaneously be determined with arbitrarily high precision.
In 1935, Einstein, Boris Podolsky, and Nathan Rosen published a claim that quantum mechanics predicts that more information about a pair of entangled particles could be observed than Heisenberg's principle allowed, which would only be possible if information were travelling instantly between the two particles. This produces a paradox which came to be known as the "EPR paradox" after the three authors. It arises if any effect felt in one location is not the result of a cause that occurred in its past light cone, relative to its location. This action at a distance seems to violate causality, by allowing information between the two locations to travel faster than the speed of light. However, it is a common misconception to think that any information can be shared between two observers faster than the speed of light using entangled particles; the hypothetical information transfer here is between the particles. See no-communication theorem for further explanation.
Based on this, the authors concluded that the quantum wave function does not provide a complete description of reality. They suggested that there must be some local hidden variables at work in order to account for the behavior of entangled particles. In a theory of hidden variables, as Einstein envisaged it, the randomness and indeterminacy seen in the behavior of quantum particles would only be apparent. For example, if one knew the details of all the hidden variables associated with a particle, then one could predict both its position and momentum. The uncertainty that had been quantified by Heisenberg's principle would simply be an artifact of not having complete information about the hidden variables. Furthermore, Einstein argued that the hidden variables should obey the condition of locality: Whatever the hidden variables actually are, the behavior of the hidden variables for one particle should not be able to instantly affect the behavior of those for another particle far away. This idea, called the principle of locality, is rooted in intuition from classical physics that physical interactions do not propagate instantly across space. These ideas were the subject of ongoing debate between their proponents. In particular, Einstein himself did not approve of the way Podolsky had stated the problem in the famous EPR paper.
In 1964, John Stewart Bell proposed his famous theorem, which states that no physical theory of hidden local variables can ever reproduce all the predictions of quantum mechanics. Implicit in the theorem is the proposition that the determinism of classical physics is fundamentally incapable of describing quantum mechanics. Bell expanded on the theorem to provide what would become the conceptual foundation of the Bell test experiments.
A typical experiment involves the observation of particles, often photons, in an apparatus designed to produce entangled pairs and allow for the measurement of some characteristic of each, such as their spin. The results of the experiment could then be compared to what was predicted by local realism and those predicted by quantum mechanics.
In theory, the results could be "coincidentally" consistent with both. To address this problem, Bell proposed a mathematical description of local realism that placed a statistical limit on the likelihood of that eventuality. If the results of an experiment violate Bell's inequality, local hidden variables can be ruled out as their cause. Later researchers built on Bell's work by proposing new inequalities that serve the same purpose and refine the basic idea in one way or another. Consequently, the term "Bell inequality" can mean any one of a number of inequalities satisfied by local hidden-variables theories; in practice, many present-day experiments employ the CHSH inequality. All these inequalities, like the original devised by Bell, express the idea that assuming local realism places restrictions on the statistical results of experiments on sets of particles that have taken part in an interaction and then separated.
To date, all Bell tests have supported the theory of quantum physics, and not the hypothesis of local hidden variables. These efforts to experimentally validate violations of the Bell inequalities resulted in John Clauser, Alain Aspect, and Anton Zeilinger being awarded the 2022 Nobel Prize in Physics.
Conduct of optical Bell test experiments
In practice most actual experiments have used light, assumed to be emitted in the form of particle-like photons (produced by atomic cascade or spontaneous parametric down conversion), rather than the atoms that Bell originally had in mind. The property of interest is, in the best known experiments, the polarisation direction, though other properties can be used. Such experiments fall into two classes, depending on whether the analysers used have one or two output channels.
A typical CHSH (two-channel) experiment
The diagram shows a typical optical experiment of the two-channel kind for which Alain Aspect set a precedent in 1982. Coincidences (simultaneous detections) are recorded, the results being categorised as '++', '+−', '−+' or '−−' and corresponding counts accumulated.
Four separate subexperiments are conducted, corresponding to the four terms E(a, b) in the test statistic S (equation (2) shown below). The settings a, a′, b and b′ are generally in practice chosen to be 0, 45°, 22.5° and 67.5° respectively — the "Bell test angles" — these being the ones for which the quantum mechanical formula gives the greatest violation of the inequality.
For each selected value of a and b, the numbers of coincidences in each category (N++, N−−, N+− and N−+) are recorded. The experimental estimate for E(a, b) is then calculated as:
Once all four E’s have been estimated, an experimental estimate of the test statistic
can be found. If S is numerically greater than 2 it has infringed the CHSH inequality. The experiment is declared to have supported the QM prediction and ruled out all local hidden-variable theories.
A strong assumption has had to be made, however, to justify use of expression (2), namely, that the sample of detected pairs is representative of the pairs emitted by the source. Denial of this assumption is called the fair sampling loophole.
A typical CH74 (single-channel) experiment
Prior to 1982 all actual Bell tests used "single-channel" polarisers and variations on an inequality designed for this setup. The latter is described in Clauser, Horne, Shimony and Holt's much-cited 1969 article as being the one suitable for practical use. As with the CHSH test, there are four subexperiments in which each polariser takes one of two possible settings, but in addition there are other subexperiments in which one or other polariser or both are absent. Counts are taken as before and used to estimate the test statistic.
where the symbol ∞ indicates absence of a polariser.
If S exceeds 0 then the experiment is declared to have infringed the CH inequality and hence to have refuted local hidden-variables. This inequality is known as CH inequality instead of CHSH as it was also derived in a 1974 article by Clauser and Horne more rigorously and under weaker assumptions.
Experimental assumptions
In addition to the theoretical assumptions, there are practical ones. There may, for example, be a number of "accidental coincidences" in addition to those of interest. It is assumed that no bias is introduced by subtracting their estimated number before calculating S, but that this is true is not considered by some to be obvious. There may be synchronisation problems — ambiguity in recognising pairs because in practice they will not be detected at exactly the same time.
Nevertheless, despite all the deficiencies of the actual experiments, one striking fact emerges: the results are, to a very good approximation, what quantum mechanics predicts. If imperfect experiments give us such excellent overlap with quantum predictions, most working quantum physicists would agree with John Bell in expecting that, when a perfect Bell test is done, the Bell inequalities will still be violated. This attitude has led to the emergence of a new sub-field of physics known as quantum information theory. One of the main achievements of this new branch of physics is showing that violation of Bell's inequalities leads to the possibility of a secure information transfer, which utilizes the so-called quantum cryptography (involving entangled states of pairs of particles).
Notable experiments
Over the past half century, a great number of Bell test experiments have been conducted. The experiments are commonly interpreted to rule out local hidden-variable theories, and in 2015 an experiment was performed that is not subject to either the locality loophole or the detection loophole (Hensen et al.). An experiment free of the locality loophole is one where for each separate measurement and in each wing of the experiment, a new setting is chosen and the measurement completed before signals could communicate the settings from one wing of the experiment to the other. An experiment free of the detection loophole is one where close to 100% of the successful measurement outcomes in one wing of the experiment are paired with a successful measurement in the other wing. This percentage is called the efficiency of the experiment. Advancements in technology have led to a great variety of methods to test Bell-type inequalities.
Some of the best known and recent experiments include:
Kasday, Ullman and Wu (1970)
Leonard Ralph Kasday, Jack R. Ullman and Chien-Shiung Wu carried out the first experimental Bell test, using photon pairs produced by positronium decay and analyzed by Compton scattering. The experiment observed photon polarization correlations consistent with quantum predictions and inconsistent with local realistic models that obey the known polarization dependence of Compton scattering. Due to the low polarization selectivity of Compton scattering, the results did not violate a Bell inequality.
Freedman and Clauser (1972)
Stuart J. Freedman and John Clauser carried out the first Bell test that observed a Bell inequality violation, using Freedman's inequality, a variant on the CH74 inequality.
Aspect et al. (1982)
Alain Aspect and his team at Orsay, Paris, conducted three Bell tests using calcium cascade sources. The first and last used the CH74 inequality. The second was the first application of the CHSH inequality. The third (and most famous) was arranged such that the choice between the two settings on each side was made during the flight of the photons (as originally suggested by John Bell).
Tittel et al. (1998)
The Geneva 1998 Bell test experiments showed that distance did not destroy the "entanglement". Light was sent in fibre optic cables over distances of several kilometers before it was analysed. As with almost all Bell tests since about 1985, a "parametric down-conversion" (PDC) source was used.
Weihs et al. (1998): experiment under "strict Einstein locality" conditions
In 1998 Gregor Weihs and a team at Innsbruck, led by Anton Zeilinger, conducted an experiment that closed the "locality" loophole, improving on Aspect's of 1982. The choice of detector was made using a quantum process to ensure that it was random. This test violated the CHSH inequality by over 30 standard deviations, the coincidence curves agreeing with those predicted by quantum theory.
Pan et al. (2000) experiment on the GHZ state
This is the first of new Bell-type experiments on more than two particles; this one uses the so-called GHZ state of three particles.
Rowe et al. (2001): the first to close the detection loophole
The detection loophole was first closed in an experiment with two entangled trapped ions, carried out in the ion storage group of David Wineland at the National Institute of Standards and Technology in Boulder. The experiment had detection efficiencies well over 90%.
Go et al. (Belle collaboration): Observation of Bell inequality violation in B mesons
Using semileptonic B0 decays of Υ(4S) at Belle experiment, a clear violation of Bell Inequality in particle-antiparticle correlation is observed.
Gröblacher et al. (2007) test of Leggett-type non-local realist theories
A specific class of non-local theories suggested by Anthony Leggett is ruled out. Based on this, the authors conclude that any possible non-local hidden-variable theory consistent with quantum mechanics must be highly counterintuitive.
Salart et al. (2008): separation in a Bell Test
This experiment filled a loophole by providing an 18 km separation between detectors, which is sufficient to allow the completion of the quantum state measurements before any information could have traveled between the two detectors.
Ansmann et al. (2009): overcoming the detection loophole in solid state
This was the first experiment testing Bell inequalities with solid-state qubits (superconducting Josephson phase qubits were used). This experiment surmounted the detection loophole using a pair of superconducting qubits in an entangled state. However, the experiment still suffered from the locality loophole because the qubits were only separated by a few millimeters.
Giustina et al. (2013), Larsson et al (2014): overcoming the detection loophole for photons
The detection loophole for photons has been closed for the first time by Marissa Giustina, using highly efficient detectors. This makes photons the first system for which all of the main loopholes have been closed, albeit in different experiments.
Christensen et al. (2013): overcoming the detection loophole for photons
The Christensen et al. (2013) experiment is similar to that of Giustina et al. Giustina et al. did just four long runs with constant measurement settings (one for each of the four pairs of settings). The experiment was not pulsed so that formation of "pairs" from the two records of measurement results (Alice and Bob) had to be done after the experiment which in fact exposes the experiment to the coincidence loophole. This led to a reanalysis of the experimental data in a way which removed the coincidence loophole, and fortunately the new analysis still showed a violation of the appropriate CHSH or CH inequality. On the other hand, the Christensen et al. experiment was pulsed and measurement settings were frequently reset in a random way, though only once every 1000 particle pairs, not every time.
Hensen et al., Giustina et al., Shalm et al. (2015): "loophole-free" Bell tests
In 2015 the first three significant-loophole-free Bell-tests were published within three months by independent groups in Delft, Vienna and Boulder. All three tests simultaneously addressed the detection loophole, the locality loophole, and the memory loophole. This makes them “loophole-free” in the sense that all remaining conceivable loopholes like superdeterminism require truly exotic hypotheses that might never get closed experimentally.
The first published experiment by Hensen et al. used a photonic link to entangle the electron spins of two nitrogen-vacancy defect centres in diamonds 1.3 kilometers apart and measured a violation of the CHSH inequality (S = 2.42 ± 0.20). Thereby the local-realist hypothesis could be rejected with a p-value of 0.039.
Both simultaneously published experiments by Giustina et al.
and Shalm et al.
used entangled photons to obtain a Bell inequality violation with high statistical significance (p-value ≪10−6). Notably, the experiment by Shalm et al. also combined three types of (quasi-)random number generators to determine the measurement basis choices. One of these methods, detailed in an ancillary file, is the “'Cultural' pseudorandom source” which involved using bit strings from popular media such as the Back to the Future films, Star Trek: Beyond the Final Frontier, Monty Python and the Holy Grail, and the television shows Saved by the Bell and Dr. Who.
Schmied et al. (2016): Detection of Bell correlations in a many-body system
Using a witness for Bell correlations derived from a multi-partite Bell inequality, physicists at the University of Basel were able to conclude for the first time Bell correlation in a many-body system composed by about 480 atoms in a Bose–Einstein condensate. Even though loopholes were not closed, this experiment shows the possibility of observing Bell correlations in the macroscopic regime.
Handsteiner et al. (2017): "Cosmic Bell Test" - Measurement Settings from Milky Way Stars
Physicists led by David Kaiser of the Massachusetts Institute of Technology and Anton Zeilinger of the Institute for Quantum Optics and Quantum Information and University of Vienna performed an experiment that "produced results consistent with nonlocality" by measuring starlight that had taken 600 years to travel to Earth. The experiment “represents the first experiment to dramatically limit the space-time region in which hidden variables could be relevant.”
Rosenfeld et al. (2017): "Event-Ready" Bell test with entangled atoms and closed detection and locality loopholes
Physicists at the Ludwig Maximilian University of Munich and the Max Planck Institute of Quantum Optics published results from an experiment in which they observed a Bell inequality violation using entangled spin states of two atoms with a separation distance of 398 meters in which the detection loophole, the locality loophole, and the memory loophole were closed. The violation of S = 2.221 ± 0.033 rejected local realism with a significance value of P = 1.02×10−16 when taking into account 7 months of data and 55000 events or an upper bound of P = 2.57×10−9 from a single run with 10000 events.
The BIG Bell Test Collaboration (2018): “Challenging local realism with human choices”
An international collaborative scientific effort used arbitrary human choice to define measurement settings instead of using random number generators. Assuming that human free will exists, this would close the “freedom-of-choice loophole”. Around 100,000 participants were recruited in order to provide sufficient input for the experiment to be statistically significant.
Rauch et al (2018): measurement settings from distant quasars
In 2018, an international team used light from two quasars (one whose light was generated approximately eight billion years ago and the other approximately twelve billion years ago) as the basis for their measurement settings. This experiment pushed the timeframe for when the settings could have been mutually determined to at least 7.8 billion years in the past, a substantial fraction of the superdeterministic limit (that being the creation of the universe 13.8 billion years ago).
The 2019 PBS Nova episode Einstein's Quantum Riddle documents this "cosmic Bell test" measurement, with footage of the scientific team on-site at the high-altitude Teide Observatory located in the Canary Islands.
Storz et al (2023): Loophole-free Bell inequality violation with superconducting circuits
In 2023, an international team led by the group of Andreas Wallraff at ETH Zurich demonstrated a loophole-free violation of the CHSH inequality with superconducting circuits deterministically entangled via a cryogenic link spanning a distance of 30 meters.
Loopholes
Though the series of increasingly sophisticated Bell test experiments has convinced the physics community that local hidden-variable theories are indefensible; they can never be excluded entirely. For example, the hypothesis of superdeterminism in which all experiments and outcomes (and everything else) are predetermined can never be excluded (because it is unfalsifiable).
Up to 2015, the outcome of all experiments that violate a Bell inequality could still theoretically be explained by exploiting the detection loophole and/or the locality loophole. The locality (or communication) loophole means that since in actual practice the two detections are separated by a time-like interval, the first detection may influence the second by some kind of signal. To avoid this loophole, the experimenter has to ensure that particles travel far apart before being measured, and that the measurement process is rapid. More serious is the detection (or unfair sampling) loophole, because particles are not always detected in both wings of the experiment. It can be imagined that the complete set of particles would behave randomly, but instruments only detect a subsample showing quantum correlations, by letting detection be dependent on a combination of local hidden variables and detector setting.
Experimenters had repeatedly voiced that loophole-free tests could be expected in the near future. In 2015, a loophole-free Bell violation was reported using entangled diamond spins over a distance of and corroborated by two experiments using entangled photon pairs.
The remaining possible theories that obey local realism can be further restricted by testing different spatial configurations, methods to determine the measurement settings, and recording devices. It has been suggested that using humans to generate the measurement settings and observe the outcomes provides a further test. David Kaiser of MIT told the New York Times in 2015 that a potential weakness of the "loophole-free" experiments is that the systems used to add randomness to the measurement may be predetermined in a method that was not detected in experiments.
Detection loophole
A common problem in optical Bell tests is that only a small fraction of the emitted photons are detected. It is then possible that the correlations of the detected photons are unrepresentative: although they show a violation of a Bell inequality, if all photons were detected the Bell inequality would actually be respected. This was first noted by Philip M. Pearle in 1970, who devised a local hidden variable model that faked a Bell violation by letting the photon be detected only if the measurement setting was favourable. The assumption that this does not happen, i.e., that the small sample is actually representative of the whole is called the fair sampling assumption.
To do away with this assumption it is necessary to detect a sufficiently large fraction of the photons. This is usually characterized in terms of the detection efficiency , defined as the probability that a photodetector detects a photon that arrives at it. Anupam Garg and N. David Mermin showed that when using a maximally entangled state and the CHSH inequality an efficiency of is required for a loophole-free violation. Later Philippe H. Eberhard showed that when using a partially entangled state a loophole-free violation is possible for , which is the optimal bound for the CHSH inequality. Other Bell inequalities allow for even lower bounds. For example, there exists a four-setting inequality which is violated for .
Historically, only experiments with non-optical systems have been able to reach high enough efficiencies to close this loophole, such as trapped ions, superconducting qubits, and nitrogen-vacancy centers. These experiments were not able to close the locality loophole, which is easy to do with photons. More recently, however, optical setups have managed to reach sufficiently high detection efficiencies by using superconducting photodetectors, and hybrid setups have managed to combine the high detection efficiency typical of matter systems with the ease of distributing entanglement at a distance typical of photonic systems.
Locality loophole
One of the assumptions of Bell's theorem is the one of locality, namely that the choice of setting at a measurement site does not influence the result of the other. The motivation for this assumption is the theory of relativity, that prohibits communication faster than light. For this motivation to apply to an experiment, it needs to have space-like separation between its measurements events. That is, the time that passes between the choice of measurement setting and the production of an outcome must be shorter than the time it takes for a light signal to travel between the measurement sites.
The first experiment that strived to respect this condition was Aspect's 1982 experiment. In it the settings were changed fast enough, but deterministically. The first experiment to change the settings randomly, with the choices made by a quantum random number generator, was Weihs et al.'s 1998 experiment. Scheidl et al. improved on this further in 2010 by conducting an experiment between locations separated by a distance of .
Coincidence loophole
In many experiments, especially those based on photon polarization, pairs of events in the two wings of the experiment are only identified as belonging to a single pair after the experiment is performed, by judging whether or not their detection times are close enough to one another. This generates a new possibility for a local hidden variables theory to "fake" quantum correlations: delay the detection time of each of the two particles by a larger or smaller amount depending on some relationship between hidden variables carried by the particles and the detector settings encountered at the measurement station.
The coincidence loophole can be ruled out entirely simply by working with a pre-fixed lattice of detection windows which are short enough that most pairs of events occurring in the same window do originate with the same emission and long enough that a true pair is not separated by a window boundary.
Memory loophole
In most experiments, measurements are repeatedly made at the same two locations. A local hidden variable theory could exploit the memory of past measurement settings and outcomes in order to increase the violation of a Bell inequality. Moreover, physical parameters might be varying in time. It has been shown that, provided each new pair of measurements is done with a new random pair of measurement settings, that neither memory nor time inhomogeneity have a serious effect on the experiment.
Superdeterminism
A necessary assumption to derive Bell's theorem is that the hidden variables are not correlated with the measurement settings. This assumption has been justified on the grounds that the experimenter has "free will" to choose the settings, and that such is necessary to do science in the first place. A (hypothetical) theory where the choice of measurement is determined by the system being measured is known as superdeterministic.
Many-worlds loophole
The many-worlds interpretation, also known as the Hugh Everett interpretation, is deterministic and has local dynamics, consisting of the unitary part of quantum mechanics without collapse. Bell's theorem does not apply because of an implicit assumption that measurements have a single outcome.
| Physical sciences | Quantum mechanics | Physics |
887427 | https://en.wikipedia.org/wiki/WinRAR | WinRAR | WinRAR is a trialware file archiver utility, developed by Eugene Roshal of win.rar GmbH. It can create and view archives in RAR or ZIP file formats, and unpack numerous archive file formats. To enable the user to test the integrity of archives, WinRAR embeds CRC32 or BLAKE2 checksums for each file in each archive. WinRAR supports creating encrypted, multi-part and self-extracting archives.
WinRAR is a Windows-only program. An Android application called "RAR for Android" is also available. Related programs include the command-line utilities "RAR" and "UNRAR" and versions for macOS, Linux, FreeBSD, WinCE, and MS-DOS.
Evolution
RAR/DOS started as a mix of x86 assembler and C, with the amount of assembly code decreasing over time and moving to pure C/C++ later on. The first versions of WinRAR were written in C, modern versions are using C++. RAR for Android is written as a mixture of Java and C++.
WinRAR and the RAR file format have evolved over time. Support for the archive format RAR5, using the same RAR file extension as earlier versions, was added in version 5.0; the older RAR file format has since been referred to as RAR4. WinRAR versions before 5.0 do not support RAR5 archives; only older versions of WinRAR run on older operating systems, and cannot open RAR5 archives.
The RAR5 file format - from version 7 on, referred to as "RAR" - increased the maximum dictionary size up to 64 GB, depending on the amount of available memory, with the default in version 5 increased from 4 MB to 32 MB, typically improving compression ratio.
For dictionaries larger than 4 GB, the size can be specified if it is unequal to a power of 2. Thus, there are no restrictions to the range 4, 8, 16, 32, 64, allowing 5 GB or 22 GB to be chosen at will. Archives with dictionaries larger than 4 GB can only be extracted by WinRAR 7.0 or newer.
AES encryption, when used, is in CBC mode and was increased in strength from 128- to 256-bit. Maximum path length for files in RAR and ZIP archives is increased from 2047 to 65535 characters.
Options added in v5.0 include 256-bit BLAKE2 file-hashing algorithm instead of default 32-bit CRC32, duplicate file detection, NTFS hard and symbolic links, and Quick Open record to allow large archives to be opened faster.
The RAR5 file format removed comments for each file (though archive comment still remains), authenticity verification, and specialized compression algorithms for text and multimedia files. RAR5 also changed the file name for split volumes from "archivename.rNN" to "archivename.partNN.rar".
The RAR7 file format added support for 64GB compression dictionary and improved compression ratio by adding two extra algorithms. RAR7 archives with dictionary sizes up to 4GB can be unpacked by previous versions of WinRAR (5.0 and above) given there's enough RAM.
Features
Creation of packed RAR or ZIP archives.
Unpacking of ARJ, BZIP2, CAB, GZ, ISO, JAR, LHA, RAR, TAR, UUE, XZ, Z, ZIP, ZIPX, ZST, 7z, UUE 001 (split) archives, as well as EXE files containing these archive formats
Checksum (integrity) verification for ARJ, BZIP2, CAB, GZ, BZIP2, RAR, XZ, ZIP and 7z archives
Multithreaded CPU compression and decompression
When creating RAR 7.0 archives:
Support for maximum file size of 16 EiB, about 1.8 × 1019 bytes or 18 million TB
Compression dictionary from 1 MiB to 64 GiB (it is limited to 256 MiB on 32-bit editions, although 32-bit editions still can decompress archives with 1 GiB dictionary; default size is 32 MiB)
Optional 256-bit BLAKE2 file hashing that can replace the default 32-bit CRC32 file checksum
Optional encryption using AES with a 256-bit key in CBC mode, using key derivation function based on PBKDF2 using HMAC-SHA256
Optional data redundancy is provided in the form of Reed–Solomon recovery records and recovery volumes, allowing reconstruction of damaged archives (including reconstruction of entirely missed volumes)
Optional "quick open record" to open RAR files faster
Ability to create multi-volume (split) archives
Ability to create self-extracting files (multi-volume self-extracting archives are supported; the self-extractor can execute commands, such as running a specified program before or after self-extraction)
Support for NTFS permissions, hard and symbolic links
Support for maximum path length up to 65,535 characters (stored in the UTF-8 format)
Optional archive comment (stored in the UTF-8 format)
Optional file timestamps preservation: modification, creation, last access times with high precision
Optional file deduplication
Advanced backup options, time-stamped files and previous file version retention.
License
The software is distributed as "try before you buy"; it may be used without charge for 40 days. When the period expires, the non-enterprise functionalities remain available, a move intended to discourage piracy. In China, a free-to-use personal edition has been provided officially since 2015.
Although archiving with the RAR format is proprietary, RARLAB supplies as copyrighted freeware the C++ source code of the current UnRAR unpacker, with a license allowing it to be used in any software, thus enabling others to produce software capable of unpacking, but not creating, RAR archives.
RAR for Android is free of charge. It displays advertisements; for a payment they can be disabled. A license for WinRAR does not provide ad-suppression for RAR for Android.
Security
In February 2019, a major security vulnerability in the unacev2.dll library which is used by WinRAR to decompress ACE archives was discovered. Consequently, WinRAR dropped the support for the ACE format from version 5.70.
Self-extracting archives created with versions before 5.31 (including the executable installer of WinRAR itself) are vulnerable to DLL hijacking: they may load and use DLLs named UXTheme.dll, RichEd32.dll and RichEd20.dll if they are in the same folder as the executable file.
It was widely reported that WinRAR v5.21 and earlier had a remote code execution (RCE) vulnerability which could allow a remote attacker to insert malicious code into a self-extracting executable (SFX) file being created by a user, "putting over 500 million users of the software at risk". However, examination of the claim revealed that, while the vulnerability existed, the result was merely an SFX which delivered its payload when executed; published responses dismissed the threat, one saying "If you can find suckers who will trust a .exe labelled as self-extracting archive ... then you can trick them into running your smuggled JavaScript".
WinRAR 6.23 fixes a critical security vulnerability which allowed the hacker to automatically execute malware distributed in archives under some circumstances.
History
Versions
Command line RAR and UNRAR were first released in autumn 1993.
Early development version WinRAR 1.54b was released in 1995.
3.00 (2002-05): the new RAR3 archive format is implemented. The new archives cannot be managed by older versions of WinRAR. Solid compression and WAV audio lossless compression features are added.
3.41 (2004-12): adds support for Linux .Z archives like GZIP and BZIP2. New options include storing entire file paths and restoring compressed NTFS files.
3.50 (2005-08): adds support for interface skins.
3.60 (2006-08): adds multithreaded version of the compression algorithm, which improves compression speed on systems with multiple dual-core or hyper-threading-enabled CPUs.
3.80 (2008-09): adds support for ZIP archives, which contain Unicode file names in UTF-8.
3.90 (2009-05): adds support for the x86-64 architecture. Multithreaded support is enhanced.
3.91 is the last release that supports Valencian.
3.92 is the last release that supports Serbian Cyrillic and Serbian Latin.
4.00 (2011-03): decompression is sped up by up to 30%.
4.10 (2012-01): removes all ZIP limitations now allowing unlimited number of files and archive size. WinRAR now also allows creation of multivolume ZIP files. ZIP archives now include Unicode file names.
4.20 (2012-06): compression speed in SMP mode is increased significantly, but this improvement was made at the expense of increased memory usage. ZIP compression now uses SMP as well. The default SMP mode cannot handle text; text compression is significantly worse unless additional switches are used.
5.00 (2013-09): the RAR5 archive format is implemented. RAR5 compressed archives cannot be managed by old versions of WinRAR. The RAR 5 format improves multi-core processor utilization, and adds a larger dictionary size of up to 1 GiB with 64-bit WinRAR. Special optional compression algorithms optimized for RGB bitmaps, raw audio files, Itanium executables, and plain text, which were supported by earlier versions, are supported only in the older RAR format, not RAR5. Optional optimized compression of x86 executables and delta compression (for structured table data) are supported in both file formats.
5.50 (2017-08): adds support for a master password which can be used to encrypt passwords stored in WinRAR. The default RAR format is changed to version 5. Adds support for decompressing Lzip archives; adds support for high precision file dates, longer file names and larger file sizes for TAR archives.
5.60 (2018-06): repairing of protected RAR5 archives was improved. Automatic detection of the encoding of ZIP archive comments. Recognition of GZIP files with arbitrary preceding data as an actual GZIP archive.
5.70 (2019-02): removes support for ACE archive decompression due to major security vulnerabilities in the unacev2.dll library.
6.00 (2020-12): "Ignore" and "Ignore All" options are added to read error prompt. "Ignore" allows to continue processing with already read file part only and "Ignore All" does it for all future read errors.
6.10 (2022-01): Added support for unpacking ZST archives. Maximum recovery record is increased to 1000% of protected data size.
6.11 (2022-03): Support of Gzip archives with large archive comments has been added; In command line mode, the switch -mes can also be used to suppress the password prompt and abort when adding new files to an encrypted solid archive.
6.12 (2022-05): security vulnerability is fixed in Unix RAR versions. WinRAR and Android RAR are not affected.
6.23 (2023-08): and critical security vulnerabilities are fixed in WinRAR. Unix and Android versions are not affected.
7.00 (2024-02): drops support for creating RAR 4.x format archives. Maximum path length limit increased to 65535 characters. Maximum RAR dictionary size up to 64 GB for the 64-bit version (limited by available RAM). Command line RAR filters out control character 27 from screen output for security reasons.
7.10 (2024-TBA): drops support for 32bit Windows editions. Adds support for Large Memory Pages which increases compression and decompression speed specially when using a large compression dictionary. Adds dark mode.
Operating systems support
More recent versions do not support many older operating systems. Versions supporting older operating systems may still be available, but not maintained:
RAR 2.50 is the last version to support MS-DOS and OS/2 on 16-bit x86 CPUs (8086-compatible).
RAR 3.93 is the last version to support MS-DOS and OS/2 on IA-32 CPUs (80386 equivalents and later). RAR for Pocket PC 3.93 is the last version for Windows Mobile. It supports file names longer than the MS-DOS standard of 8.3 characters, in a MS-DOS box (except under NT-based operating systems), and uses the RSX DPMI extender.
WinRAR 2.06 is the last version to support Windows 3.1, Windows NT 3.1, Windows NT 3.5, Windows NT 3.51 and Win32s.
WinRAR 3.93 is the last version to support Windows 95, Windows NT 4.0, Windows 98 and Windows Me.
WinRAR 4.11 is the last version to support Windows 2000.
WinRAR 6.02 is the last version to support Windows XP (except the console version Rar.exe).
WinRAR 7.01 is the last version to support Windows Vista (and 32-bit Windows editions).
| Technology | Office and data management | null |
887815 | https://en.wikipedia.org/wiki/Hair%20conditioner | Hair conditioner | Hair conditioner is a hair care cosmetic product used to improve the feel, texture, appearance and manageability of hair. Its main purpose is to reduce friction between strands of hair to allow smoother brushing or combing, which might otherwise cause damage to the scalp. Various other benefits are often advertised, such as hair repair, strengthening, or a reduction in split ends.
Conditioners are available in a wide range of forms, including viscous liquids, gels and creams, as well as thinner lotions and sprays. Hair conditioner is usually used after the hair has been washed with shampoo. It is applied and worked into the hair and may either be rinsed out a short time later or left in.
History
For centuries, natural oils have been used to condition human hair. A conditioner popular with men in the late Victorian era was Macassar oil, but this product was quite greasy and necessitated the pinning of a small cloth, known as an antimacassar, to the headrests of chairs and sofas to preserve the upholstery from being damaged by the oil.
Modern hair conditioner was created at the turn of the 20th century when the Edouard Pinaud company presented a product he called Brilliantine at the 1900 Exposition Universelle in Paris. His product was intended to soften men's hair, including beards and moustaches. Since the invention of Pinaud's early products, modern science has advanced the hair conditioner industry to include those made with silicone, fatty alcohols and quaternary ammonium compounds. These chemical products have the benefits of hair conditioner without feeling greasy or heavy.
Mechanism of action
The outermost layer of a hair follicle is called the cuticle and is composed largely of keratin. This is rich in cysteine groups which are mildly acidic. When the hair is washed these groups can deprotonate, giving the hair a negative charge.
The ingredients in conditioner, especially positively charged quaternary ammonium species, such as behentrimonium chloride or polymers that are known as polyquaternium-XX, where XX is an arbitrary number, can then become attached to the hair via electrostatic interactions. Once attached these compounds have several effects. Their long hydrocarbon backbone helps to lubricate the surface of each hair follicle, reducing the sensation of roughness and assisting combing. The surface coating of cationic groups means that hairs are repelled from each other electrostatically, which reduces clumping. The compounds can also act as antistatic agents, which helps to reduce frizzing.
Types
Conditioners, also called deep conditioners or hair masks, are heavy and thick, with a high content of cationic surfactants that are able to bind to the hair structure and "glue" the hair surface scales together. This type of conditioner is designed to restore hair's moisture levels and reduce breakage. These are usually applied to the hair for a longer time (30–45 minutes).
Leave-in conditioners are thinner and have different surfactants, which add only a little material to the hair to avoid weighing down the hair or causing greasiness. They are based on unsaturated fatty acid chains, which are bent, not straight. Leave-in conditioner is designed to be used in a similar way to hair oil, preventing the tangling of hair and keeping it smooth. Its use is particularly prevalent among those with naturally curly or kinky hair.
Rinse-out/rinse-through conditioners are the most common or generic on the market. Ordinary conditioners are generally applied directly after using shampoo, and manufacturers usually produce a conditioner counterpart for different types of shampoo for this purpose.
Hold conditioners, based on cationic polyelectrolyte polymers, hold the hair in a desired shape. These have a function and composition similar to diluted hair gels.
Cleansing conditioners are a newer category, typically based on a combination of amphoteric and cationic surfactants that can be used either in place of shampoo or as a pretreatment before shampooing for hair that is damaged or very curly.
Ingredients
There are several types of hair conditioner ingredients, differing in composition and functionality:
Acidifiers are acidity regulators that maintain the conditioner's pH at about 3.5. In contact with an acidic environment, the hair's somewhat scaly surface tightens up as the hydrogen bonds between the keratin molecules are strengthened.
Antistatic agents, which bind to the hair and reduce the static, can include cationic polymers such as polyquaternium-10 and guar hydroxypropyltrimonium chloride.
Detanglers modify the hair surface pH as acidifiers or coat it with polymers as glossers.
Glossers are light-reflecting chemicals that bind to the hair surface and are usually polymers, usually silicones (e.g., dimethicone or cyclopentasiloxane).
Lubricants such as fatty alcohols, panthenol, dimethicone, etc.
Moisturizers, whose role is to hold moisture in the hair, usually contain high proportions of humectants. These could also be provided by natural oils such as Prunus Amygdalus Dulcis (sweet almond) oil.
Oils (EFAs – essential fatty acids) can help dry/porous hair become more soft and pliable. The scalp produces a natural oil called sebum. Sebum naturally contains EFAs.
Preservatives protect the product from spoilage by microorganisms during the product's shelf life.
Reconstructors, usually containing hydrolyzed protein, supposedly penetrate the hair and strengthen its structure through polymer cross-linking.
Sequestrants improve function in hard water.
Sunscreen provides protection against protein degradation and color loss. Currently, benzophenone-4 and ethylhexyl methoxycinnamate are the two sunscreens most commonly used in hair products. Cinnamidopyltrimonium chloride and a few others are used to a much lesser degree. The common sunscreens used on the skin are rarely used for hair products due to their texture and weight effects.
Surfactants – approximately 97% of hair consists of a protein called keratin. The surface of keratin contains negatively charged amino acids. Hair conditioners therefore usually contain cationic surfactants, which don't wash out completely, because their hydrophilic ends strongly bind to keratin. The hydrophobic ends of the surfactant molecules then act as the new hair surface. Examples are behentrimonium chloride and cetrimonium chloride.
Thermal protectors, usually heat-absorbing polymers, shield the hair against excessive heat caused by blow-drying, curling irons, hot rollers, etc.
pH
Conditioners are frequently acidic, as low pH protonates the keratin's amino acids. The hydrogen ions give the hair a positive charge and create more hydrogen bonds among the keratin scales, giving the hair a more compact structure. Organic acids such as citric acid are usually used to maintain acidity.
| Biology and health sciences | Hygiene products | Health |
888273 | https://en.wikipedia.org/wiki/Cantilever%20bridge | Cantilever bridge | A cantilever bridge is a bridge built using structures that project horizontally into space, supported on only one end (called cantilevers). For small footbridges, the cantilevers may be simple beams; however, large cantilever bridges designed to handle road or rail traffic use trusses built from structural steel, or box girders built from prestressed concrete.
The steel truss cantilever bridge was a major engineering breakthrough when first put into practice, as it can span distances of over , and can be more easily constructed at difficult crossings by virtue of using little or no falsework.
Origins
Engineers in the 19th century understood that a bridge that was continuous across multiple supports would distribute the loads among them. This would result in lower stresses in the girder or truss and meant that longer spans could be built. Several 19th-century engineers patented continuous bridges with hinge points mid-span. The use of a hinge in the multi-span system presented the advantages of a statically determinate system and of a bridge that could handle differential settlement of the foundations. Engineers could more easily calculate the forces and stresses with a hinge in the girder.
Heinrich Gerber was one of the engineers to obtain a patent for a hinged girder (1866) and is recognized as the first to build one. The Hassfurt Bridge over the Main river in Germany with a central span of 124 feet (38 metres) was completed in 1867 and is recognized as the first modern cantilever bridge.
The High Bridge of Kentucky by C. Shaler Smith (1877), the Niagara Cantilever Bridge by Charles Conrad Schneider (1883) and the Poughkeepsie Bridge by John Francis O'Rourke and Pomeroy P. Dickinson (1889) were all important early uses of the cantilever design. The Kentucky River Bridge spanned a gorge that was 275 feet (84 metres) deep and took full advantage of the fact that falsework, or temporary support, is not needed for the main span of a cantilever bridge.
The Forth Bridge is a notable example of an early cantilever bridge. This bridge held the record for longest span in the world for twenty-nine years until it was surpassed by the Quebec Bridge. The engineers responsible for the bridge, Sir Benjamin Baker and Sir John Fowler, demonstrated the structural principles of the suspended span cantilever by sitting in chairs and supporting their colleague, Kaichi Watanabe, in between them, using just their arms and wooden poles. The suspended span, where Watanabe sits, is in the center. The wooden poles resist the compression of the lower chord, while the outstretched arms support the tension of the upper chord. The placement of the brick counterweights demonstrates the action of the outer foundations.
Function
A simple cantilever span is formed by two cantilever arms extending from opposite sides of an obstacle to be crossed, meeting at the center. In a common variant, the suspended span, the cantilever arms do not meet in the center; instead, they support a central truss bridge which rests on the ends of the cantilever arms. The suspended span may be built off-site and lifted into place, or constructed in place using special travelling supports.
A common way to construct steel truss and prestressed concrete cantilever spans is to counterbalance each cantilever arm with another cantilever arm projecting the opposite direction, forming a balanced cantilever; when they attach to a solid foundation, the counterbalancing arms are called anchor arms. Thus, in a bridge built on two foundation piers, there are four cantilever arms: two which span the obstacle, and two anchor arms that extend away from the obstacle. Because of the need for more strength at the balanced cantilever's supports, the bridge superstructure often takes the form of towers above the foundation piers. The Commodore Barry Bridge is an example of this type of cantilever bridge.
Steel truss cantilevers support loads by tension of the upper members and compression of the lower ones. Commonly, the structure distributes the tension via the anchor arms to the outermost supports, while the compression is carried to the foundations beneath the central towers. Many truss cantilever bridges use pinned joints and are therefore statically determinate with no members carrying mixed loads.
Prestressed concrete balanced cantilever bridges are often built using segmental construction.
Construction methods
Some steel arch bridges (such as the Navajo Bridge) are built using pure cantilever spans from each side, with neither falsework below nor temporary supporting towers and cables above. These are then joined with a pin, usually after forcing the union point apart, and when jacks are removed and the bridge decking is added the bridge becomes a truss arch bridge. Such unsupported construction is only possible where appropriate rock is available to support the tension in the upper chord of the span during construction, usually limiting this method to the spanning of narrow canyons.
List by length
World's longest cantilever bridges (by longest span):
Examples
| Technology | Transport infrastructure | null |
888711 | https://en.wikipedia.org/wiki/Mathematical%20statistics | Mathematical statistics | Mathematical statistics is the application of probability theory and other mathematical concepts to statistics, as opposed to techniques for collecting statistical data. Specific mathematical techniques that are commonly used in statistics include mathematical analysis, linear algebra, stochastic analysis, differential equations, and measure theory.
Introduction
Statistical data collection is concerned with the planning of studies, especially with the design of randomized experiments and with the planning of surveys using random sampling. The initial analysis of the data often follows the study protocol specified prior to the study being conducted. The data from a study can also be analyzed to consider secondary hypotheses inspired by the initial results, or to suggest new studies. A secondary analysis of the data from a planned study uses tools from data analysis, and the process of doing this is mathematical statistics.
Data analysis is divided into:
descriptive statistics – the part of statistics that describes data, i.e. summarises the data and their typical properties.
inferential statistics – the part of statistics that draws conclusions from data (using some model for the data): For example, inferential statistics involves selecting a model for the data, checking whether the data fulfill the conditions of a particular model, and with quantifying the involved uncertainty (e.g. using confidence intervals).
While the tools of data analysis work best on data from randomized studies, they are also applied to other kinds of data. For example, from natural experiments and observational studies, in which case the inference is dependent on the model chosen by the statistician, and so subjective.
Topics
The following are some of the important topics in mathematical statistics:
Probability distributions
A probability distribution is a function that assigns a probability to each measurable subset of the possible outcomes of a random experiment, survey, or procedure of statistical inference. Examples are found in experiments whose sample space is non-numerical, where the distribution would be a categorical distribution; experiments whose sample space is encoded by discrete random variables, where the distribution can be specified by a probability mass function; and experiments with sample spaces encoded by continuous random variables, where the distribution can be specified by a probability density function. More complex experiments, such as those involving stochastic processes defined in continuous time, may demand the use of more general probability measures.
A probability distribution can either be univariate or multivariate. A univariate distribution gives the probabilities of a single random variable taking on various alternative values; a multivariate distribution (a joint probability distribution) gives the probabilities of a random vector—a set of two or more random variables—taking on various combinations of values. Important and commonly encountered univariate probability distributions include the binomial distribution, the hypergeometric distribution, and the normal distribution. The multivariate normal distribution is a commonly encountered multivariate distribution.
Special distributions
Normal distribution, the most common continuous distribution
Bernoulli distribution, for the outcome of a single Bernoulli trial (e.g. success/failure, yes/no)
Binomial distribution, for the number of "positive occurrences" (e.g. successes, yes votes, etc.) given a fixed total number of independent occurrences
Negative binomial distribution, for binomial-type observations but where the quantity of interest is the number of failures before a given number of successes occurs
Geometric distribution, for binomial-type observations but where the quantity of interest is the number of failures before the first success; a special case of the negative binomial distribution, where the number of successes is one.
Discrete uniform distribution, for a finite set of values (e.g. the outcome of a fair die)
Continuous uniform distribution, for continuously distributed values
Poisson distribution, for the number of occurrences of a Poisson-type event in a given period of time
Exponential distribution, for the time before the next Poisson-type event occurs
Gamma distribution, for the time before the next k Poisson-type events occur
Chi-squared distribution, the distribution of a sum of squared standard normal variables; useful e.g. for inference regarding the sample variance of normally distributed samples (see chi-squared test)
Student's t distribution, the distribution of the ratio of a standard normal variable and the square root of a scaled chi squared variable; useful for inference regarding the mean of normally distributed samples with unknown variance (see Student's t-test)
Beta distribution, for a single probability (real number between 0 and 1); conjugate to the Bernoulli distribution and binomial distribution
Statistical inference
Statistical inference is the process of drawing conclusions from data that are subject to random variation, for example, observational errors or sampling variation. Initial requirements of such a system of procedures for inference and induction are that the system should produce reasonable answers when applied to well-defined situations and that it should be general enough to be applied across a range of situations. Inferential statistics are used to test hypotheses and make estimations using sample data. Whereas descriptive statistics describe a sample, inferential statistics infer predictions about a larger population that the sample represents.
The outcome of statistical inference may be an answer to the question "what should be done next?", where this might be a decision about making further experiments or surveys, or about drawing a conclusion before implementing some organizational or governmental policy.
For the most part, statistical inference makes propositions about populations, using data drawn from the population of interest via some form of random sampling. More generally, data about a random process is obtained from its observed behavior during a finite period of time. Given a parameter or hypothesis about which one wishes to make inference, statistical inference most often uses:
a statistical model of the random process that is supposed to generate the data, which is known when randomization has been used, and
a particular realization of the random process; i.e., a set of data.
Regression
In statistics, regression analysis is a statistical process for estimating the relationships among variables. It includes many ways for modeling and analyzing several variables, when the focus is on the relationship between a dependent variable and one or more independent variables. More specifically, regression analysis helps one understand how the typical value of the dependent variable (or 'criterion variable') changes when any one of the independent variables is varied, while the other independent variables are held fixed. Most commonly, regression analysis estimates the conditional expectation of the dependent variable given the independent variables – that is, the average value of the dependent variable when the independent variables are fixed. Less commonly, the focus is on a quantile, or other location parameter of the conditional distribution of the dependent variable given the independent variables. In all cases, the estimation target is a function of the independent variables called the regression function. In regression analysis, it is also of interest to characterize the variation of the dependent variable around the regression function which can be described by a probability distribution.
Many techniques for carrying out regression analysis have been developed. Familiar methods, such as linear regression, are parametric, in that the regression function is defined in terms of a finite number of unknown parameters that are estimated from the data (e.g. using ordinary least squares). Nonparametric regression refers to techniques that allow the regression function to lie in a specified set of functions, which may be infinite-dimensional.
Nonparametric statistics
Nonparametric statistics are values calculated from data in a way that is not based on parameterized families of probability distributions. They include both descriptive and inferential statistics. The typical parameters are the expectations, variance, etc. Unlike parametric statistics, nonparametric statistics make no assumptions about the probability distributions of the variables being assessed.
Non-parametric methods are widely used for studying populations that take on a ranked order (such as movie reviews receiving one to four stars). The use of non-parametric methods may be necessary when data have a ranking but no clear numerical interpretation, such as when assessing preferences. In terms of levels of measurement, non-parametric methods result in "ordinal" data.
As non-parametric methods make fewer assumptions, their applicability is much wider than the corresponding parametric methods. In particular, they may be applied in situations where less is known about the application in question. Also, due to the reliance on fewer assumptions, non-parametric methods are more robust.
One drawback of non-parametric methods is that since they do not rely on assumptions, they are generally less powerful than their parametric counterparts. Low power non-parametric tests are problematic because a common use of these methods is for when a sample has a low sample size. Many parametric methods are proven to be the most powerful tests through methods such as the Neyman–Pearson lemma and the Likelihood-ratio test.
Another justification for the use of non-parametric methods is simplicity. In certain cases, even when the use of parametric methods is justified, non-parametric methods may be easier to use. Due both to this simplicity and to their greater robustness, non-parametric methods are seen by some statisticians as leaving less room for improper use and misunderstanding.
Statistics, mathematics, and mathematical statistics
Mathematical statistics is a key subset of the discipline of statistics. Statistical theorists study and improve statistical procedures with mathematics, and statistical research often raises mathematical questions.
Mathematicians and statisticians like Gauss, Laplace, and C. S. Peirce used decision theory with probability distributions and loss functions (or utility functions). The decision-theoretic approach to statistical inference was reinvigorated by Abraham Wald and his successors and makes extensive use of scientific computing, analysis, and optimization; for the design of experiments, statisticians use algebra and combinatorics. But while statistical practice often relies on probability and decision theory, their application can be controversial
| Mathematics | Statistics | null |
889183 | https://en.wikipedia.org/wiki/Inclinometer | Inclinometer | An inclinometer or clinometer is an instrument used for measuring angles of slope, elevation, or depression of an object with respect to gravity's direction. It is also known as a tilt indicator, tilt sensor, tilt meter, slope alert, slope gauge, gradient meter, gradiometer, level gauge, level meter, declinometer, and pitch & roll indicator. Clinometers measure both inclines and declines using three different units of measure: degrees, percentage points, and topos. The astrolabe is an example of an inclinometer that was used for celestial navigation and location of astronomical objects from ancient times to the Renaissance.
A tilt sensor can measure the tilting in often two axes of a reference plane in two axes.
In contrast, a full motion would use at least three axes and often additional sensors. One way to measure tilt angle with reference to the earth's ground plane, is to use an accelerometer. Typical applications can be found in the industry and in game controllers. In aircraft, the "ball" in turn coordinators or turn and bank indicators is sometimes referred to as an inclinometer.
History
Inclinometers include examples such as Well's in-clinometer, the essential parts of which are a flat side, or base, on which it stands, and a hollow disc just half filled with some heavy liquid. The glass face of the disc is surrounded by a graduated scale that marks the angle at which the surface of the liquid stands, with reference to the flat base. The zero line is parallel to the base, and when the liquid stands on that line, the flat side is horizontal; the 90 degree is perpendicular to the base, and when the liquid stands on that line, the flat side is perpendicular or plumb. Intervening angles are marked, and, with the aid of simple conversion tables, the instrument indicates the rate of fall per set distance of horizontal measurement, and set distance of the sloping line.
Al-Biruni, a Persian polymath, once wanted to measure the height of the sun. He lacked the necessary equipment to measure this height. He was forced to create a calibrated arc on the back of a counting board, which he then used as a makeshift quadrant with the help of a plumb line. He determined the location's latitude using the measurements taken with this rudimentary tool. This quadrant was most likely an inclinometer based on the quarter-circle panel.
The Abney level is a handheld surveying instrument developed in the 1870s that includes a sighting tube and inclinometer, arranged so that the surveyor may align the sighting tube (and its crosshair) with the reflection of the bubble in the spirit level of the inclinometer when the line of sight is at the angle set on the inclinometer.
One of the more famous inclinometer installations was on the panel of the Ryan NYP "The Spirit of St. Louis"—in 1927 Charles Lindbergh chose the lightweight Rieker Inc P-1057 Degree Inclinometer to give him climb and descent angle information.
Uses
Hand-held clinometers are used for a variety of surveying and measurement tasks. In land surveying and mapping, a clinometer can provide a rapid measurement of the slope of a geographic feature, or used for cave survey. In prospecting for minerals, clinometers are used to measure the strike and dip of geologic formations. In forestry, tree height measurement can be done with a clinometer using standardized methods including triangulation. Major artillery guns may have an associated clinometer used to facilitate aiming of shells over long distances.
Permanently-installed tiltmeters are emplaced at major earthworks such as dams to monitor the long-term stability of the structure.
Factors which influence the use of inclinometers
(Overall accuracy varies depending on the type of tilt sensor (or inclinometer) and technology used)
Gravity
Temperature (drift), zero offset, linearity, vibration, shock, cross-axis sensitivity, acceleration/deceleration.
A clear line of sight between the user and the measured point is needed.
A well defined object is required to obtain the maximum precision.
The angle measurement precision and accuracy is limited to slightly better than one arcsec.
Accuracy
Certain highly sensitive electronic inclinometer sensors can achieve an output resolution to 0.0001°; depending on the technology and angle range, it may be limited to 0.001°. An inclinometer sensor's true or absolute accuracy (which is the combined total error), however, is a combination of initial sets of sensor zero offset and sensitivity, sensor linearity, hysteresis, repeatability, and the temperature drifts of zero and sensitivity—electronic inclinometers accuracy can typically range from ±0.01–2° depending on the sensor and situation. Typically in room ambient conditions the accuracy is limited to the sensor linearity specification.
Sensor technology
Tilt sensors and inclinometers generate an artificial horizon and measure angular tilt with respect to this horizon. They are used in cameras, aircraft flight controls, automobile security systems, and speciality switches and are also used for platform leveling, boom angle indication, and in other applications requiring measurement of tilt.
Important specifications to consider for tilt sensors and inclinometers are the tilt angle range and the number of axes. The axes are usually, but not always, orthogonal. The tilt angle range is the range of desired linear output.
Common implementations of tilt sensors and inclinometers are accelerometer, Liquid Capacitive, electrolytic, gas bubble in liquid, and pendulum.
Tilt sensor technology has also been implemented in video games. Yoshi's Universal Gravitation and Kirby Tilt 'n' Tumble are both built around a tilt sensor mechanism, which is built into the cartridge. The PlayStation 3 and Wii game controllers also use tilt as a means to play video games.
Inclinometers are also used in civil engineering, for example, to measure the inclination of land to be built upon.
Some inclinometers provide an electronic interface based on CAN (Controller Area Network). In addition, those inclinometers may support the standardized CANopen profile (CiA 410). In this case, these inclinometers are compatible and partly interchangeable.
Two-axis digital inclinometer
Traditional spirit levels and pendulum-based electronic leveling instruments are usually constrained by only single-axis and narrow tilt measurement range. However, most precision leveling, angle measurement, alignment and surface flatness profiling tasks essentially involve a two-dimensional surface plane angle rather than two independent orthogonal single-axis objects. Two-axis inclinometers that are built with MEMS tilt sensors provides simultaneous two-dimensional angle readings of a surface plane tangent to earth datum.
Typical advantages of using two-axis MEMS inclinometers over conventional single-axis "bubble" or mechanical leveling instruments may include:
Simultaneous measurement of two-dimensional (X-Y plane) tilt angles (i.e. pitch & roll), can eliminate tedious swapping back-and-forth experienced when using a single-axis level, for example to adjust machine footings to attain a precise leveling position.
Digital compensation and precise calibration for non-linearity, for example for operating temperature variation, resulting in higher accuracy over a wider measurement range.
The accelerometer sensors may generate numerical data in the form of vibration profiles to enable a machine installer to track and assess alignment quality in real-time and verify a structure's positional stability by comparing leveling profiles before and after it is set up.
Inclinometer with gyroscope
As inclinometers measure the angle of an object with respect to the force of gravity, external accelerations like rapid motions, vibrations or shocks will introduce errors in the tilt measurements. To overcome this problem, it is possible to use a gyroscope in addition to an accelerometer. Any of the abovementioned accelerations have a huge impact on the accelerometer, but a limited effect on the measured rotation rates of the gyroscope. An algorithm can combine both signals to get the best value out of each sensor. This way it is possible to separate the actual tilt angle from the errors introduced by external accelerations.
Applications
Inclinometers are used for:
Determining latitude using Polaris (in the Northern Hemisphere) or the two stars of the constellation Crux (in the Southern Hemisphere).
Determining the angle of the Earth's magnetic field with respect to the horizontal plane.
Showing a deviation from the true vertical or horizontal.
Surveying, to measure an angle of inclination or elevation.
Alerting an equipment operator that it may tip over.
Measuring angles of elevation, slope, or incline, e.g. of an embankment.
Measuring slight differences in slopes, particularly for geophysics. Such inclinometers are, for instance, used for monitoring volcanoes, or for measuring the depth and rate of landslide movement.
Measuring movements in walls or the ground in civil engineering projects.
Determining the dip of beds or strata, or the slope of an embankment or cutting; a kind of plumb level.
Some automotive safety systems.
Indicating pitch and roll of vehicles, nautical craft, and aircraft. See turn coordinator and slip indicator.
Monitoring the boom angle of cranes and material handlers.
Measuring the "look angle" of a satellite antenna towards a satellite.
Adjusting a solar panel to the optimal angle to maximize its output.
Measuring the slope angle of a tape or chain during distance measurement.
Measuring the height of a building, tree, or other feature using a vertical angle and a distance (determined by taping or pacing), using trigonometry.
Measuring the angle of drilling in well logging.
Measuring the list of a ship in still water and the roll in rough water.
Measuring steepness of a ski slope.
Measuring the orientation of planes and lineations in rocks, in combination with a compass, in structural geology.
Measuring range of motion in the joints of the body
Measuring the inclination angle of the pelvis. Numerous neck and back measurements require the simultaneous use of two inclinometers.
it measures the angle of elevation, and ultimately computing the altitudes of, many things otherwise inaccessible for direct measurement.
Measuring and fine tuning the angle of line array speaker hangs. Confirmation of the angle achieved via use of a laser built into the remote inclinometer.
Setting correct orientation of solar panels while installing
Setting firing angle of a cannon or gun (determines projectile range)
Electronic games
Help prevent unsafe working conditions.
The USDA Forest Service uses tilt sensors (or inclinometers) to measure tree height in its Forest Inventory and Analysis program.
Logistics and transport also use tilt indicators, it is a specific system for single use. They are attached to the products during the shipping process.
Games
Nintendo used tilt sensor technology in five games for its Game Boy series of hand-held game systems. The tilt sensor allows players to control aspects of the game by twisting the game system. Games that use this feature:
Yoshi's Universal Gravitation (Game Boy Advance)
WarioWare: Twisted! (Game Boy Advance)(not released in Europe)
Koro Koro Puzzle Happy Panechu! (Game Boy Advance)(Japan only)
Kirby Tilt 'n' Tumble (Game Boy Color)(not released in Europe)
Command Master (Game Boy Color)(Japan only)
Tilt sensors can also be found in game controllers such as the Microsoft Sidewinder Freestyle Pro and Sony's PlayStation 3 controller.
However, unlike these other controllers in which the tilt sensor serves as a supplement to normal control methods, it serves as one of the central features of Nintendo's Wii Remote and the Nunchuk attachment. Along with accelerometers, the tilt sensors are a primary method of control in most Wii games.
It is now being used in many different aspects, instead of just games like motocrossing and flight simulators. It can be used for sport gaming, first-person shooter, and other odd uses such as in WarioWare: Smooth Moves
Another example is a virtual version of a wooden maze with obstacles in which you have to maneuver a ball by tilting the maze. A homebrew tilt sensor interface was made for the Palm (PDA).
| Technology | Surveying tools | null |
889856 | https://en.wikipedia.org/wiki/Superheated%20steam | Superheated steam | Superheated steam is steam at a temperature higher than its vaporization point at the absolute pressure where the temperature is measured.
Superheated steam can therefore cool (lose internal energy) by some amount, resulting in a lowering of its temperature without changing state (i.e., condensing) from a gas to a mixture of saturated vapor and liquid. If unsaturated steam (a mixture which contains both water vapor and liquid water droplets) is heated at constant pressure, its temperature will also remain constant as the vapor quality (think dryness, or percent saturated vapor) increases towards 100%, and becomes dry (i.e., no saturated liquid) saturated steam. Continued heat input will then "super" heat the dry saturated steam. This will occur if saturated steam contacts a surface with a higher temperature.
Superheated steam and liquid water cannot coexist under thermodynamic equilibrium, as any additional heat simply evaporates more water and the steam will become saturated steam. However, this restriction may be violated temporarily in dynamic (non-equilibrium) situations. To produce superheated steam in a power plant or for processes (such as drying paper) the saturated steam drawn from a boiler is passed through a separate heating device (a superheater) which transfers additional heat to the steam by contact or by radiation.
Superheated steam is not suitable for sterilization. This is because the superheated steam is dry. Dry steam must reach much higher temperatures and the materials exposed for a longer time period to have the same effectiveness; or equal F0 kill value. Superheated steam is also not useful for heating; while it has more energy and can do more work than saturated steam, its heat content is much less useful. This is because superheated steam has the same heat transfer coefficient of air, making it an insulator - a poor conductor of heat. Saturated steam has a much higher wall heat transfer coefficient.
Slightly superheated steam may be used for antimicrobial disinfection of biofilms on hard surfaces.
Superheated steam's greatest value lies in its tremendous internal energy that can be used for kinetic reaction through mechanical expansion against turbine blades and reciprocating pistons, that produces rotary motion of a shaft. The value of superheated steam in these applications is its ability to release tremendous quantities of internal energy yet remain above the condensation temperature of water vapor; at the pressures at which reaction turbines and reciprocating piston engines operate.
Of prime importance in these applications is the fact that water vapor containing entrained liquid droplets is generally incompressible at those pressures. In a reciprocating engine or turbine, if steam doing work cools to a temperature at which liquid droplets form, then the water droplets entrained in the fluid flow will strike the mechanical parts with enough force to bend, crack or fracture them. Superheating and pressure reduction through expansion ensures that the steam flow remains as a compressible gas throughout its passage through a turbine or an engine, preventing damage of the internal moving parts.
Saturated steam
Saturated steam is steam that is in equilibrium with heated water at the same pressure, i.e., it has not been heated above the boiling point for its pressure. This is in contrast to superheated steam, in which the steam (vapor) has been separated from the water droplets then additional heat has been added.
These condensation droplets are a cause of damage to steam turbine blades, the reason why such turbines rely on a supply of dry, superheated steam.
Dry steam is saturated steam that has been very slightly superheated. This is not sufficient to change its energy appreciably, but is a sufficient rise in temperature to avoid condensation problems, given the average loss in temperature across the steam supply circuit. Towards the end of the 19th century, when superheating was still a less-than-certain technology, such steam-drying gave the condensation-avoiding benefits of superheating without requiring the sophisticated boiler or lubrication techniques of full superheating.
By contrast, water vapor that includes water droplets is described as wet steam. If wet steam is heated further, the droplets evaporate, and at a high enough temperature (which depends on the pressure) all of the water evaporates, the system is in vapor–liquid equilibrium, and it becomes saturated steam.
Saturated steam is advantageous in heat transfer due to the high latent heat of vaporization. It is a very efficient mode of heat transfer. In layman's terms, saturated steam is at its dew point at the corresponding temperature and pressure. The typical latent heat of vaporization (or condensation) is for saturated steam at atmospheric pressure.
Uses
Steam engine
Superheated steam was widely used in main line steam locomotives. Saturated steam has three main disadvantages in a steam engine: it contains small droplets of water which have to be periodically drained from the cylinders; being precisely at the boiling point of water for the boiler pressure in use, it inevitably condenses to some extent in the steam pipes and cylinders outside the boiler, causing a disproportionate loss of steam volume as it does so; and it places a heavy demand on the boiler.
Superheating the steam dries it effectively, raises its temperature to a point where condensation is much less likely and increases its volume significantly. Added together, these factors increase the power and economy of the locomotive. The main disadvantages are the added complexity and cost of the superheater tubing and the adverse effect that the "dry" steam has on lubrication of moving components such as the steam valves. Shunting locomotives did not generally use superheating.
The normal arrangement involved taking steam after the regulator valve and passing it through long superheater tubes inside specially large firetubes of the boiler. The superheater tubes had a reverse ("torpedo") bend at the firebox end so that the steam had to pass the length of the boiler at least twice, picking up heat as it did so.
Processing
Other potential uses of superheated steam include: drying, cleaning, layering, reaction engineering, epoxy drying and film use where saturated to highly superheated steam is required at one atmospheric pressure or at high pressure. Ideal for steam drying, steam oxidation and chemical processing. Uses are in surface technologies, cleaning technologies, steam drying, catalysis, chemical reaction processing, surface drying technologies, curing technologies, energy systems and nanotechnologies.
The application of superheated steam for sanitation of dry food processing plant environment has been reported.
Superheated steam is not usually used in a heat exchanger due to low heat transfer co-efficient. In refining and hydrocarbon industries superheated steam is mainly used for stripping and cleaning purposes.
Pest control
Steam has been used for soil steaming since the 1890s. Steam is induced into the soil which causes almost all organic material to deteriorate (the term "sterilization" is used, but it is not strictly correct since all micro-organism are not necessarily killed). Soil steaming is an effective alternative to many chemicals in agriculture, and is used widely by greenhouse growers. Wet steam is primarily used in this process, but if soil temperatures above the boiling point of water are required, superheated steam must be used.
| Physical sciences | Phase transitions | Physics |
889940 | https://en.wikipedia.org/wiki/Trivial%20name | Trivial name | In chemistry, a trivial name is a non-systematic name for a chemical substance. That is, the name is not recognized according to the rules of any formal system of chemical nomenclature such as IUPAC inorganic or IUPAC organic nomenclature. A trivial name is not a formal name and is usually a common name.
Generally, trivial names are not useful in describing the essential properties of the thing being named. Properties such as the molecular structure of a chemical compound are not indicated. And, in some cases, trivial names can be ambiguous or will carry different meanings in different industries or in different geographic regions (for example, a trivial name such as white metal can mean various things). Trivial names are simpler. As a result, a limited number of trivial chemical names are retained names, an accepted part of the nomenclature.
Trivial names often arise in the common language; they may come from historic usages in, for example, alchemy. Many trivial names pre-date the institution of formal naming conventions. Names can be based on a property of the chemical, including appearance (color, taste or smell), consistency, and crystal structure; a place where it was found or where the discoverer comes from; the name of a scientist; a mythological figure; an astronomical body; the shape of the molecule; and even fictional figures. All elements that have been isolated have trivial names.
Definitions
In scientific documents, international treaties, patents and legal definitions, names for chemicals are needed that
identify them unambiguously. This need is satisfied by systematic names. One such system, established by the International Union of Pure and Applied Chemistry (IUPAC), was established in 1950. Other systems have been developed by the American Chemical Society, the International Organization for Standardization, and the World Health Organization. However, chemists still use many names that are not systematic because they are traditional or because they are more convenient than the systematic names. These are called trivial names. The word "trivial", often used in a pejorative sense, was intended to mean "commonplace".
In addition to trivial names, chemists have constructed semi-trivial names by appending a standard symbol to a trivial stem. Some trivial and semi-trivial names are so widely used that they have been officially adopted by IUPAC; these are known as retained names.
Pesticide common names
The common names used for pesticides did not become commonplace through repeated informal usage, The names are granted by ISO committee (TC81), who approve the common name according to ISO1750.
Elements
Traditional names of elements are trivial, some originating in alchemy. IUPAC has accepted these names, but has also defined systematic names of elements that have not yet been prepared. It has adopted a procedure by which the scientists who are credited with preparing an element can propose a new name. Once the IUPAC has accepted such a (trivial) name, it replaces the systematic name.
Origins
Nine elements were known by the Middle Ages: gold, silver, tin, mercury, copper, lead, iron, sulfur, and carbon. Mercury was named after the planet, but its symbol was derived from the Latin hydrargyrum, which itself comes from the Greek υδράργυρος, meaning liquid silver; mercury is also known as quicksilver in English. The symbols for the other eight are derived from their Latin names.
Systematic nomenclature began after Louis-Bernard Guyton de Morveau stated the need for "a constant method of denomination, which helps the intelligence and relieves the memory". The resulting system was popularized by Antoine Lavoisier's publication of Méthode de nomenclature chimique (Method of Chemical Nomenclature) in 1787. Lavoisier proposed that elements be named after their properties. For the next 125 years, most chemists followed this suggestion, using Greek and Latin roots to compose the names; for example, hydrogen ("water-producing"), oxygen ("acid-producing"), nitrogen ("soda-producing"), bromine ("stink"), and argon were based on Greek roots, while the names of iodine and chlorine were derived from the Greek words for their characteristic colors. Indium, rubidium, and thallium were similarly named for the colors of particular lines in their emission spectra. Iridium, which forms compounds of many different colors, takes its name from iris, the Latin for "rainbow". The noble gases have all been named for their origin or properties. Helium comes from the Greek helios, meaning "Sun" because it was first detected as a line in the spectrum of the Sun (it is not known why the suffix -ium, which is used for metals, was chosen). The other noble gases are neon ("new"), argon ("slow, lazy"), krypton ("hidden"), xenon ("stranger"), and radon ("from radium").
Many more elements have been given names that have little or nothing to do with their properties. Elements have been named for celestial bodies (helium, selenium, tellurium, for the Sun, Moon, and Earth; cerium and palladium for Ceres and Pallas, two asteroids). They have been named for mythological figures, including Titans in general (titanium) and Prometheus in particular (promethium); Roman and Greek gods (uranium, neptunium, and plutonium) and their descendants (tantalum for Tantalus, a son of Zeus, and niobium for Niobe, a daughter of Tantalus); and Norse deities (vanadium for the goddess Vanadis and thorium for the god Thor).
Some elements were named for aspects of the history of their discovery. In particular, technetium and promethium were so named because the first samples detected were artificially synthesised; neither of the two has any isotope sufficiently stable to occur in nature on Earth in significant quantities. The connection to the Titan Prometheus was that he had been fabled to have stolen fire from the gods for mankind.
Discoverers of some elements named them after their home country or city. Marie Curie named polonium after Poland; ruthenium, gallium, germanium, and lutetium were based on the Latin names for Russia, France, Germany, and Paris. Other elements are named after the place where they were discovered. Four elements — terbium, erbium, ytterbium, and yttrium — were named after the Swedish village Ytterby, where ores containing them were extracted. Other elements named after places are magnesium (after Magnesia), strontium, scandium, europium, thulium (after an old Roman name for an unidentified northern region), holmium, copper (derived from Cyprus, where it was mined in the Roman era), hafnium, rhenium, americium, berkelium, californium, and darmstadtium.
For the elements up to 92 (uranium), naming elements after people was discouraged. The two exceptions are indirect, the elements being named after minerals that were themselves named after people. These were gadolinium (found in gadolinite, named after the Finnish chemist Johan Gadolin) and samarium (the mineral samarskite was named after a Russian mining engineer, Vasili Samarsky-Bykhovets). Among the transuranium elements, this restriction was relaxed; there followed curium (after the Curies), einsteinium (Albert Einstein), fermium (Enrico Fermi), mendelevium (Dmitri Mendeleev), nobelium (Alfred Nobel) and lawrencium (Ernest Lawrence).
Relation to IUPAC standards
IUPAC has established international standards for naming elements. The first scientist or laboratory to isolate an element has the right to propose a name; after a review process, a final decision is made by the IUPAC Council. In keeping with tradition, names can be based on a mythological concept or character, astronomical object, mineral, place, property of the element or scientist. For those elements that have not yet been discovered, IUPAC has established a systematic name system. The names combine syllables that represent the digits of the atomic number, followed by "-ium". For example, "unununium" is element 111 ("un" being the syllable for 1). However, once the element has been found, the systematic name is replaced by a trivial one, in this case roentgenium.
The IUPAC names for elements are intended for use in the official languages. At the time of the first edition of the IUPAC Red Book (which contains the rules for inorganic compounds), those languages were English and French; now English is the sole official language. However, other languages still have their own names for elements. The chemical symbol for tungsten, W, is based on the German name , which is found in wolframite and comes from the German for "wolf's foam", how the mineral was known to Saxon miners. The name tungsten means "heavy stone", a description of scheelite, another mineral in which tungsten is found. Russian names for hydrogen, oxygen and carbon are vodorod, kislorod and uglerod (generating water, acid and coal respectively). The German names for hydrogen, oxygen, and nitrogen are (water substance), (acid substance), and (smothering substance).
The corresponding Chinese names are qīngqì (light gas), yǎngqì (nourishing gas), and dànqì (diluting gas). A method for translating chemical names into Chinese was developed by John Fryer and Xu Shou in 1871. Where traditional names were well established, they kept them; otherwise, a single character was created.
Inorganic chemistry
Early terminology for compound chemicals followed similar rules to the naming of elements. The names could be based on the appearance of the substance, including all five senses. In addition, chemicals were named after the consistency, crystalline form, a person or place, its putative medical properties or method of preparation.
Salt (sodium chloride) is soluble and is used to enhance the taste of food. Substances with similar properties came to be known as salts, in particular Epsom salt (magnesium sulfate, found in a bitter saline spring in the English town of Epsom). Ammonium (with the little-used systematic name azanium) was first extracted from sal ammoniac, meaning "salt of Amun". Ancient Romans noticed crystals of it in Egyptian temples devoted to the god Amun; the crystals had condensed from the smoke of burning camel dung. Lead acetate was called sugar of lead. However, other names like sugar of lead (lead(II) acetate), butter of antimony (antimony trichloride), oil of vitriol (sulfuric acid), and cream of tartar (potassium bitartrate) borrowed their language from the kitchen. Many more names were based on color; for example, hematite, orpiment, and verdigris come from words meaning "blood-like stone", "gold pigment", and "green of Greece".
Some names are based on their use. Lime is a general name for materials combining calcium with carbonates, oxides or hydroxides; the name comes from a root "sticking or adhering"; its earliest use was as mortar for construction.
Water has several systematic names, including oxidane (the IUPAC name), hydrogen oxide, and dihydrogen monoxide (DHMO). The latter was the basis of the dihydrogen monoxide hoax, a document that was circulated warning readers of the dangers of the chemical (for example, it is fatal if inhaled).
Organic chemistry
In organic chemistry, some trivial names derive from a notable property of the thing being named. For instance, lecithin, the common name for phosphatidylcholine, was originally isolated from egg yolk. The word is coined from the Greek λέκιθος (lékithos) for yolk.
Many trivial names continue to be used because their sanctioned equivalents are considered too cumbersome for everyday use. For example, "tartaric acid", a compound found in wine, has a systematic name of 2,3-dihydroxybutanedioic acid. The pigment β-Carotene has an IUPAC name of 1,3,3-trimethyl-2-[(1E,3E,5E,7E,9E,11E,13E,15E,17E)-3,7,12,16-tetramethyl-18-(2,6,6-trimethylcyclohexen-1-yl)octadeca-1,3,5,7,9,11,13,15,17-nonaenyl]cyclohexene. However, the trivial name can be potentially confusing. Based on its name, one might come to the conclusion that the molecule theobromine contains one or more bromine atoms. In reality it is an alkaloid similar in structure to caffeine.
Shape-based
Several organic molecules have semitrivial names where the suffixes -ane (for an alkane) or -ene (for an alkene) are added to a name based on the shape of the molecule. Some are pictured below. Other examples include barrelene (shaped like a barrel), fenestrane (having a window-pane motif), ladderane (a ladder shape), olympiadane (having a shape with the same topology as the Olympic rings) and quadratic acid (also known as squaric acid).
Based on fiction
The bohemic acid complex is a mixture of chemicals obtained through fermentation of a species of actinobacteria. In 1977 the components were isolated and have been found useful as antitumor agents and anthracycline antibiotics. The authors named the complex (and one of its components, bohemamine) after the opera La bohème by Puccini, and the remaining components were named after characters in the opera: alcindoromycin (Alcindoro), collinemycin (Colline), marcellomycin (Marcello), mimimycin (Mimi), musettamycin (Musetta), rudolphomycin (Rodolfo) and schaunardimycin (Schaunard). However, the relationships between the characters do not correctly reflect the chemical relationships.
A research lab at Lepetit Pharmaceuticals, led by Piero Sensi, was fond of coining nicknames for chemicals that they discovered, later converting them to a form more acceptable for publication. The antibiotic rifampicin was named after a French movie, Rififi, about a jewel heist. They nicknamed another antibiotic "Mata Hari" before changing the name to matamycin.
| Physical sciences | Nomenclature | Chemistry |
890069 | https://en.wikipedia.org/wiki/Thymol | Thymol | Thymol (also known as 2-isopropyl-5-methylphenol, IPMP), , is a natural monoterpenoid phenol derivative of p-Cymene, isomeric with carvacrol. It occurs naturally in the oil of thyme, and it is extracted from Thymus vulgaris (common thyme), ajwain, and various other plants as a white crystalline substance of a pleasant aromatic odor and strong antiseptic properties. Thymol also provides the distinctive, strong flavor of the culinary herb thyme, also produced from T. vulgaris. Thymol is only slightly soluble in water at neutral pH, but it is extremely soluble in alcohols and other organic solvents. It is also soluble in strongly alkaline aqueous solutions due to deprotonation of the phenol. Its dissociation constant (pKa) is . Thymol absorbs maximum UV radiation at 274 nm.
Chemical synthesis
Thymol is produced by the alkylation of m-cresol and propene:
A predicted method of biosynthesis of thymol in thyme and oregano begins with the cyclization of geranyl diphosphate by TvTPS2 to γ-terpinene. Oxidation by a cytochrome P450 in the CYP71D subfamily creates a dienol intermediate, which is then converted into a ketone by short-chain dehydrogenase. Lastly, keto-enol tautomerization gives thymol.
History
Ancient Egyptians used thyme for embalming. The ancient Greeks used it in their baths and burned it as incense in their temples, believing it was a source of courage. The spread of thyme throughout Europe was thought to be due to the Romans, as they used it to purify their rooms and to "give an aromatic flavour to cheese and liqueurs". In the European Middle Ages, the herb was placed beneath pillows to aid sleep and ward off nightmares. In this period, women also often gave knights and warriors gifts that included thyme leaves, because it was believed to bring courage to the bearer. Thyme was also used as incense and placed on coffins during funerals, because it was supposed to ensure passage into the next life.
The bee balms Monarda fistulosa and Monarda didyma, North American wildflowers, are natural sources of thymol. The Blackfoot Native Americans recognized these plants' strong antiseptic action and used poultices of the plants for skin infections and minor wounds. A tisane made from them was also used to treat mouth and throat infections caused by dental caries and gingivitis.
Thymol was first isolated by German chemist Caspar Neumann in 1719. In 1853, French chemist Alexandre Lallemand (1816-1886) named thymol and determined its empirical formula. Antiseptic properties of thymol were discovered in 1875, and it was first synthesized by Swedish chemist Oskar Widman (1852-1930) in 1882.
Extraction
The conventional method of extracting is hydro-distillation (HD), but can also be extracted with solvent-free microwave extraction (SFME). In 30 minutes, SFME yields similar amounts of thymol with more oxygenated compounds than 4.5 hours of hydro-distillation at atmospheric pressures without the need for solvent.
Uses
Thymol during the 1910s was the treatment of choice for hookworm infection in the United States. People of the Middle East continue to use za'atar, a delicacy made with large amounts of thyme, to reduce and eliminate internal parasites. It is also used as a preservative in halothane, an anaesthetic, and as an antiseptic in mouthwash. When used to reduce plaque and gingivitis, thymol has been found to be more effective when used in combination with chlorhexidine than when used purely by itself. Thymol is also the active antiseptic ingredient in some toothpastes, such as Johnson & Johnson's Euthymol. Thymol has been used to successfully control varroa mites and prevent fermentation and the growth of mold in bee colonies. Thymol is also used as a rapidly degrading, non-persisting pesticides such as insecticides and fungicides which are leveraged in plant care products, where its environmentally friendly, rapid degradation ensures it doesn’t leave persistent residues while effectively controlling pests and fungal issues. Thymol can also be used as a medical disinfectant and general purpose disinfectant. Thymol is also used in the production of menthol through the hydrogenation of the aromatic ring.
List of plants that contain thymol
Illicium verum
Euphrasia rostkoviana
Lagoecia cuminoides
Monarda didyma
Monarda fistulosa
Mosla chinensis
Ocimum gratissimum L.
Origanum compactum
Origanum dictamnus
Origanum onites
Origanum vulgare
Satureja hortensis
Satureja thymbra
Thymus glandulosus
Thymus hyemalis
Thymus serpyllum
Thymus praecox
Thymus vulgaris
Thymus zygis
Trachyspermum ammi
Toxicology and environmental impacts
In 2009, the U.S. Environmental Protection Agency (EPA) reviewed the research literature on the toxicology and environmental impact of thymol and concluded that "thymol has minimal potential toxicity and poses minimal risk".
Environmental breakdown and use as a pesticide
Studies have shown that hydrocarbon monoterpenes and thymol in particular degrade rapidly (DT50 16 days in water, 5 days in soil) in the environment and are, thus, low risks because of rapid dissipation and low bound residues, supporting the use of thymol as a pesticide agent that offers a safe alternative to other more persistent chemical pesticides that can be dispersed in runoff and produce subsequent contamination. Though, there has been recent research into sustained released systems for botanically derived pesticides, such as using natural polysaccharides which would be biodegradable and biocompatible.
Compendial status
British Pharmacopoeia
Japanese Pharmacopoeia
| Physical sciences | Terpenes and terpenoids | Chemistry |
890149 | https://en.wikipedia.org/wiki/Alpine%20ibex | Alpine ibex | The Alpine ibex (Capra ibex), also known as the steinbock, is a European species of goat that lives in the Alps. It is one of ten species in the genus Capra and its closest living relative is the Iberian ibex. The Alpine ibex is a sexually dimorphic species; males are larger and carry longer horns than females. Its coat is brownish-grey. Alpine ibexes tend to live in steep, rough terrain and open alpine meadows. They can be found at elevations as high as and their sharp hooves allow them to scale the steep slopes and cliffs of their mountainous habitat.
Alpine ibexes primarily feed on grass and are active throughout the year. Although they are social animals, adult males and females segregate for most of the year, coming together only to mate. During the breeding season, males use their long horns to fight for access to females. Ibexes have few predators but may succumb to parasites and diseases.
By the 19th century, the Alpine ibex had been extirpated from most of its range and it went through a population bottleneck of fewer than 100 individuals during its near-extinction event, leading to very low genetic diversity across populations. The species has been successfully reintroduced to parts of its historical range. All individuals living today descend from the stock in Gran Paradiso National Park, Italy. , the IUCN lists the species as being of least concern.
Taxonomy
Carl Linnaeus first described the Alpine ibex in 1758. It is classified in the genus Capra with nine other species of goat. is Latin for while the species name is translated from Latin as and is possibly derived from an earlier Alpine language.
Fossils of the genus Tossunnoria are found in late Miocene deposits in China; these fossils appear to have been transitional between goats and their ancestors. The genus Capra may have originated in Central Asia and spread to Europe, the Caucasus, and East Africa from the Pliocene and into the Pleistocene. Mitochondrial and Y chromosome evidence show hybridisation of species in this lineage. Fossils of the Alpine ibex dating from the last glacial period during the late Pleistocene have been found in France and Italy. The Alpine Ibex and the Iberian ibex (C. pyrenaica) probably evolved from the extinct Pleistocene species Capra camburgensis, whose fossils have been found in Germany. The Alpine ibex appears to have been larger during the Pleistocene than in the modern day.
In the 20th century, the Nubian (C. nubiana), walia (C. walie), and Siberian ibex (C. sibirica) were considered to be subspecies of the Alpine ibex; populations in the Alps were given the trinomial of C. i. ibex. Genetic evidence from 2006 has supported the status of these Ibexes as separate species.
The following cladogram of seven Capra species is based on 2022 mitochondrial evidence:
Appearance
Alpine ibexes are sexually dimorphic. Males grow to a height of at the withers with a body length of and weigh . Females are much smaller and have a shoulder height of , a body length of , and weigh .
The Alpine ibex is a stocky animal with a tough neck and robust legs with short metapodials. Compared with most other wild goats, the species has a wide, shortened snout. Adaptations for climbing include sharp, highly separated hooves and a rubbery callus under the front feet. Both male and female Alpine ibexes have large, backwards-curving horns with an elliptical cross-section and a trilateral-shaped core. Transverse ridges on the front surface of the horns mark an otherwise flat surface. At , the horns of males are substantially longer than those of females, which reach only in length.
The species has brownish-grey hair over most of its body but lighter in colour on the belly with dark markings on the chin and throat. The hair on the chest region is nearly black and there are stripes along the dorsal (back) surface. The Alpine ibex is duller-coloured than other members of its genus. As with other goats, only males have a beard. Ibexes moult in spring, when their thick winter coat consisting of woolly underfur is replaced with a short, thin summer coat. Their winter coat grows back in the autumn. As in other members of Capra, the Alpine ibex has glands near the eyes, groin, and feet but none on the face.
Distribution and habitat
The Alpine ibex is native to the Alps of central Europe; its range includes France, Switzerland, Liechtenstein, Italy, Germany, and Austria. Fossils of the species have been found as far south as Greece, where it became locally extinct about 7,500 years ago due to human predation. Between the 16th and 18th centuries, the species disappeared from much of its range due to hunting, leaving by the 19th century one surviving population in and around Gran Paradiso, Italy. The species has since been reintroduced into parts of its former range, as well as new areas such as Slovenia and Bulgaria.
The Alpine Ibex is an excellent climber; it occupies steep, rough terrain at elevations of . It prefers to live in open areas, but when there is little snow and depending on population density, adult males may gather in larch and mixed larch-spruce woodland. Outside the breeding season, the sexes live in separate habitats. Females are more likely to be found on steep slopes while males prefer more level ground. Males inhabit lowland meadows during the spring, when fresh grass appears, and climb to alpine meadows during the summer. In early winter, both males and females move to steep, rocky slopes to avoid dense buildups of snow. Alpine ibexes prefer slopes of 30–45°, and take refuge in small caves and overhangs.
Behaviour and ecology
The Alpine ibex is strictly herbivorous; its diet consists mostly of grass, which is preferred all year; during the summer, ibexes supplement their diet with herbs, while during autumn and winter they also eat dwarf shrubs and conifer shoots. The most-commonly eaten grass genera are Agrostis, Avena, Calamagrostis, Festuca, Phleum, Poa, Sesleria, and Trisetum. In the spring, animals of both sexes spend about the same amount of time feeding during the day, while in summer, females, particularly those that are lactating, eat more than males. High temperatures cause heat stress in large adult males, reducing their feeding time, but they may avoid this problem by feeding at night.
In Gran Paradiso, home ranges of the Alpine Ibex can exceed and in reintroduced populations, home ranges may approach . Home-range size depends on the availability of resources and the time of year. Home ranges tend to be largest during summer and autumn, smallest in winter, and intermediate in spring. Females' home ranges are usually smaller than those of males. Ibexes do not hibernate during the winter; they take shelter on cold winter nights and bask in the mornings. They also reduce their heart rate and metabolism.
The Alpine ibex may compete for resources with chamois and red deer; the presence of these species may force the ibex to occupy higher elevations. The Alpine ibex's climbing ability is such that it has been observed scaling the 57-degree slopes of the Cingino Dam in Piedmont, Italy, where it licks salts. Only females and kids, which are lighter and have shorter legs than adult males, will climb the steep dam. Kids have been observed at , ascending in a zig-zag path while descending in straight paths.
Social behaviour
The Alpine ibex is a social species but it tends to live in groups that are based on sex and age. For most of the year, adult males group separately from females, and older males live separately from young males. Female groups consist of 5–10 members and male groups usually have 2–16 members but sometimes have more than 50. Dependent kids live with their mothers in female groups. Segregation between the sexes is a gradual process; males younger than nine years may still associate with female groups. Adult males, particularly older males, are more likely to be found alone than females. Social spacing tends to be looser in the summer, when there is more room to feed. Ibexes have stable social connections; they consistently regroup with the same individuals when ecological conditions force them together. Female groups tend to be more stable than male groups.
Adult males and females gather together in December and January, the breeding season, then separate again in April and May. Among males, a dominance hierarchy based on size, age, and horn length exists. Hierarchies are established outside the breeding season, allowing males to focus more on mating and less on fighting. Males use their horns for combat; they bash rivals' sides or clash head-to-head often by rearing then clashing downwards.
Alpine ibexes communicate mainly through short, sharp whistles that serve mostly as alarm calls and may occur singularly or in succession with short gaps. Females and their young communicate by bleating.
Reproduction and growth
The mating season begins in December and typically lasts for around six weeks. During this time, male herds break up into smaller groups and search for females. The rut takes place in two phases; in the first phase, males interact with females as a group and in the second phase, one male separates from his group to follow a female in oestrus. Dominant males between nine and twelve years old follow a female and guard her from rivals while subordinate, younger males between two and six years old try to sneak past the tending male when he is distracted. If the female flees, both dominant and subordinate males will try to follow her. During courtship, the male stretches the neck, flicks the tongue, curls the upper lip, urinates, and sniffs the female. After copulation, the male rejoins his group and restarts the first phase of the rut. Environmental conditions can affect courtship in the species; for example, snow can limit the males' ability to follow females and mate with them.
The female is in oestrus for around 20 days and gestation averages around five months, and typically results in the birth of one or sometimes two kids. Females give birth away from their social groups on rocky slopes that are relatively safe from predators. After a few days, the kids can move on their own. Mothers and kids gather into nursery groups, where young are nursed for up to five months. Nursery groups can also include non-lactating females. Alpine ibexes reach sexual maturity in 18 months but continue to grow until females are five to six and males are nine to eleven years old.
The horns grow throughout life. Young are born without horns, which become visible as tiny tips at one month and reach in the second month. In males, the horns grow at about per year for the first five-and-a-half years, slowing to half that rate once the animal reaches 10 years of age. The slowing of horn growth in males coincides with aging. The age of an ibex can be determined by annual growth rings in the horns, which stop growing in winter.
Mortality and health
Male Alpine Ibexes live for around 16 years while females live for around 20 years. The species has a high adult survival rate compared with other herbivores around its size. In one study, all kids reached two years of age and the majority of adults lived for 13 years, although most 13-year-old males did not reach the age of 15. Alpine ibexes have a low rate of predation; their mountain habitat keeps them safe from predators like wolves, though golden eagles may prey on young. In Gran Paradiso, causes of death are old age, lack of food, and disease. They are also killed by avalanches.
Alpine ibexes may suffer necrosis and fibrosis caused by the bacteria Brucella melitensis, and foot rot caused by Dichelobacter nodosus. Infections from Mycoplasma conjunctivae damage the eye via keratoconjunctivitis and can lead to death rates of up to 30%. Ibexes can host gastrointestinal parasites such as coccidia, strongyles, Teladorsagia circumcincta, and Marshallagi amarshalli as well as lungworms, mainly Muellerius capillaris. Several individuals have died from heart diseases, including arteriosclerosis, cardiac fibrosis, sarcosporidiosis, and valvular heart disease.
Conservation
During the Middle Ages, the Alpine ibex ranged throughout the Alpine region of Europe. Starting in the early 16th century, the overall population declined due almost entirely to hunting by humans, especially with the introduction of firearms. By the 19th century, only around 100 individuals remained in and around Gran Paradiso in north-west Italy and on the Italian-French border. In 1821, the Government of Piedmont banned hunting of the Alpine ibex and in 1854, Victor Emmanuel II declared Gran Paradiso a royal hunting reserve. In 1920, his grandson Victor Emmanuel III of Italy donated the land to the state of Italy and it was established as a national park. By 1933, the Alpine ibex population reached 4,000 but subsequent mismanagement by the Fascist government caused it to drop to around 400 by 1945. Their protection improved after the war and by 2005, there were 4,000 in the national park. In the late 20th century, the Gran Paradiso population was used for reintroductions into other parts of Italy.
Starting in 1902, several Alpine ibexes from Gran Paradiso were taken into captive facilities in Switzerland for selective breeding and reintroduction into the wild. Until 1948, translocated founder animals were captive-bred. Afterwards, there were reintroductions of wild-born specimens from established populations in Piz Albris, Le Pleureur, and Augstmatthorn. These gave rise to the populations in France and Austria. Alpine ibexes also recolonised areas on their own. The Alpine ibex population reached 3,020 in 1914, 20,000 in 1991, and 55,297 in 2015, and by 1975, the species occupied much of its medieval range. In the 1890s, ibexes were introduced to Slovenia despite the lack of evidence of their presence there following the last glacial period. In 1980, ibexes were translocated to Bulgaria.
Between 2015 and 2017, there were around 9,000 ibexes in 30 colonies in France, over 17,800 individuals and 30 colonies in Switzerland, over 16,400 ibexes in 67 colonies in Italy, around 9,000 in 27 colonies in Austria, around 500 in five colonies in Germany, and almost 280 ibexes and four colonies in Slovenia. As of 2020, the IUCN considers the Alpine ibex to be of Least Concern with a stable population trend. It was given a recovery score of 79%, making it "moderately depleted". While the species would likely have gone extinct without conservation efforts in the 19th and 20th centuries, as of 2021, it has a low conservation dependence. According to the IUCN, without current protections, the population decline of the species would be minimal. Some countries allow limited hunting.
Having gone through a genetic bottleneck, the Alpine ibex population has low genetic diversity and is at risk of inbreeding depression. A 2020 analysis found highly deleterious mutations were lost in these new populations but they had also gained mildly deleterious ones. The genetic purity of the species may be threatened by hybridisation with domestic goats, which have been allowed to roam in the Alpine Ibex's habitat. The genetic bottleneck of populations may increase vulnerability to infectious diseases because their immune system has low major histocompatibility complex diversity. In the Bornes Massif region of the French Alps, management actions, including a test-and-cull program to control outbreaks, effectively reduced Brucella infection prevalence in adult females from 51% in 2013 to 21% in 2018, and active infections also significantly declined.
Cultural significance
The Alpine ibex is called the steinbock, which originated from the Old High German word steinboc, literally "stone buck". Several European names for the animal developed from this, including the French bouquetin and the Italian stambecco. The Alpine ibex is one of many animals depicted in the art of the Late Pleistocene-era Magdalenian culture in Western Europe. Local people used Ibexes for traditional medicine; the horn material was used to counter cramps, poisoning, and hysteria, while the blood was thought to prevent stones from developing in the bladder. The species' value as a source of medicine led to its near extinction. Since its recovery, the Alpine ibex has been seen as a resilient symbol of the mountain range. The species is depicted on the coat of arms of the Swiss canton of Grisons.
| Biology and health sciences | Bovidae | Animals |
890232 | https://en.wikipedia.org/wiki/Hair%20cell | Hair cell | Hair cells are the sensory receptors of both the auditory system and the vestibular system in the ears of all vertebrates, and in the lateral line organ of fishes. Through mechanotransduction, hair cells detect movement in their environment.
In mammals, the auditory hair cells are located within the spiral organ of Corti on the thin basilar membrane in the cochlea of the inner ear. They derive their name from the tufts of stereocilia called hair bundles that protrude from the apical surface of the cell into the fluid-filled cochlear duct. The stereocilia number from fifty to a hundred in each cell while being tightly packed together and decrease in size the further away they are located from the kinocilium.
Mammalian cochlear hair cells are of two anatomically and functionally distinct types, known as outer, and inner hair cells. Damage to these hair cells results in decreased hearing sensitivity, and because the inner ear hair cells cannot regenerate, this damage is permanent. Damage to hair cells can cause damage to the vestibular system and therefore cause difficulties in balancing. However, other vertebrates, such as the frequently studied zebrafish, and birds have hair cells that can regenerate.
The human cochlea contains on the order of 3,500 inner hair cells and 12,000 outer hair cells at birth.
The outer hair cells mechanically amplify low-level sound that enters the cochlea. The amplification may be powered by the movement of their hair bundles, or by an electrically driven motility of their cell bodies. This so-called somatic electromotility amplifies sound in all tetrapods. It is affected by the closing mechanism of the mechanical sensory ion channels at the tips of the hair bundles.
The inner hair cells transform the sound vibrations in the fluids of the cochlea into electrical signals that are then relayed via the auditory nerve to the auditory brainstem and to the auditory cortex.
Inner hair cells – from sound to nerve signal
The deflection of the hair-cell stereocilia opens mechanically gated ion channels that allow any small, positively charged ions (primarily potassium and calcium) to enter the cell. Unlike many other electrically active cells, the hair cell itself does not fire an action potential. Instead, the influx of positive ions from the endolymph in the scala media depolarizes the cell, resulting in a receptor potential. This receptor potential opens voltage gated calcium channels; calcium ions then enter the cell and trigger the release of neurotransmitters at the basal end of the cell. The neurotransmitters diffuse across the narrow space between the hair cell and a nerve terminal, where they then bind to receptors and thus trigger action potentials in the nerve. In this way, the mechanical sound signal is converted into an electrical nerve signal. Repolarization of hair cells is done in a special manner. The perilymph in the scala tympani has a very low concentration of positive ions. The electrochemical gradient makes the positive ions flow through channels to the perilymph.
Hair cells chronically leak Ca2+. This leakage causes a tonic release of neurotransmitter to the synapses. It is thought that this tonic release is what allows the hair cells to respond so quickly in response to mechanical stimuli. The quickness of the hair cell response may also be due to the fact that it can increase the amount of neurotransmitter release in response to a change of as little as 100 μV in membrane potential.
Hair cells are also able to distinguish tone frequencies through one of two methods. The first method, found only in non-mammals, uses electrical resonance in the basolateral membrane of the hair cell. The electrical resonance for this method appears as a damped oscillation of membrane potential responding to an applied current pulse. The second method uses tonotopic differences of the basilar membrane. This difference comes from the different locations of the hair cells. Hair cells that have high-frequency resonance are located at the basal end while hair cells that have significantly lower frequency resonance are found at the apical end of the epithelium.
Outer hair cells – acoustical pre-amplifiers
In mammalian outer hair cells, the varying receptor potential is converted to active vibrations of the cell body. This mechanical response to electrical signals is termed somatic electromotility;
it drives variations in the cell's length, synchronized to the incoming sound signal, and provides mechanical amplification by feedback to the traveling wave.
Outer hair cells are found only in mammals. While hearing sensitivity of mammals is similar to that of other classes of vertebrates, without functioning outer hair cells, the sensitivity decreases by approximately 50 dB. Outer hair cells extend the hearing range to about 200 kHz in some marine mammals. They have also improved frequency selectivity (frequency discrimination), which is of particular benefit for humans, because it enabled sophisticated speech and music. Outer hair cells are functional even after cellular stores of ATP are depleted.
The effect of this system is to nonlinearly amplify quiet sounds more than large ones so that a wide range of sound pressures can be reduced to a much smaller range of hair displacements. This property of amplification is called the cochlear amplifier.
The molecular biology of hair cells has seen considerable progress in recent years, with the identification of the motor protein (prestin) that underlies somatic electromotility in the outer hair cells. Prestin's function has been shown to be dependent on chloride channel signaling and that it is compromised by the common marine pesticide tributyltin. Because this class of pollutant bioconcentrates up the food chain, the effect is pronounced in top marine predators such as orcas and toothed whales.
Hair cell signal adaptation
Calcium ion influx plays an important role for the hair cells to adapt to the amplification of the signal. This allows humans to ignore constant sounds that are no longer new and allow us to be acute to other changes in our surrounding. The key adaptation mechanism comes from a motor protein myosin-1c that allows slow adaptation, provides tension to sensitize transduction channels, and also participate in signal transduction apparatus. More recent research now shows that the calcium-sensitive binding of calmodulin to myosin-1c could actually modulate the interaction of the adaptation motor with other components of the transduction apparatus as well.
Fast Adaptation: During fast adaptation, Ca2+ ions that enter a stereocilium through an open MET channel bind rapidly to a site on or near the channel and induce channel closure. When channels close, tension increases in the tip link, pulling the bundle in the opposite direction. Fast adaptation is more prominent in sound and auditory detecting hair cells, rather in vestibular cells.
Slow Adaption: The dominating model suggests that slow adaptation occurs when myosin-1c slides down the stereocilium in response to elevated tension during bundle displacement. The resultant decreased tension in the tip link permits the bundle to move farther in the opposite direction. As tension decreases, channels close, producing the decline in transduction current. Slow adaptation is most prominent in vestibular hair cells that sense spatial movement and less in cochlear hair cells that detect auditory signals.
Neural connection
Neurons of the auditory or vestibulocochlear nerve (the eighth cranial nerve) innervate cochlear and vestibular hair cells. The neurotransmitter released by hair cells that stimulates the terminal neurites of peripheral axons of the afferent (towards the brain) neurons is thought to be glutamate. At the presynaptic juncture, there is a distinct presynaptic dense body or ribbon. This dense body is surrounded by synaptic vesicles and is thought to aid in the fast release of neurotransmitter.
Nerve fiber innervation is much denser for inner hair cells than for outer hair cells. A single inner hair cell is innervated by numerous nerve fibers, whereas a single nerve fiber innervates many outer hair cells. Inner hair cell nerve fibers are also very heavily myelinated, which is in contrast to the unmyelinated outer hair cell nerve fibers. The region of the basilar membrane supplying the inputs to a particular afferent nerve fibre can be considered to be its receptive field.
Efferent projections from the brain to the cochlea also play a role in the perception of sound. Efferent synapses occur on outer hair cells and on afferent axons under inner hair cells. The presynaptic terminal bouton is filled with vesicles containing acetylcholine and a neuropeptide called calcitonin gene-related peptide. The effects of these compounds vary; in some hair cells the acetylcholine hyperpolarizes the cell, which reduces the sensitivity of the cochlea locally.
Regrowth
Research on the regrowth of cochlear cells may lead to medical treatments that restore hearing. Unlike birds and fish, humans and other mammals are generally incapable of regrowing the cells of the inner ear that convert sound into neural signals when those cells are damaged by age or disease. Researchers are making progress in gene therapy and stem-cell therapy that may allow the damaged cells to be regenerated. Because hair cells of auditory and vestibular systems in birds and fish have been found to regenerate, their ability has been studied at length. In addition, lateral line hair cells, which have a mechanotransduction function and are found in anamniotes, have been shown to regrow in species such as the zebrafish.
Researchers have identified a mammalian gene that normally acts as a molecular switch to block the regrowth of cochlear hair cells in adults. The Rb1 gene encodes the retinoblastoma protein, which is a tumor suppressor. Rb stops cells from dividing by encouraging their exit from the cell cycle. Not only do hair cells in a culture dish regenerate when the Rb1 gene is deleted, but mice bred to be missing the gene grow more hair cells than control mice that have the gene. Additionally, the sonic hedgehog protein has been shown to block activity of the retinoblastoma protein, thereby inducing cell cycle re-entry and the regrowth of new cells.
Several Notch signaling pathway inhibitors, including the gamma secretase inhibitor LY3056480, are being studied for their potential ability to regenerate hair cells in the cochlea.
TBX2 (T-box transcription factor 2) has been shown to be a master regulator in the differentiation of inner and outer hair cells. This discovery has allowed researchers to direct hair cells to develop into either inner or outer hair cells, which could help in replacing hair cells that have died and prevent or reverse hearing loss.
The cell cycle inhibitor p27kip1 (CDKN1B) has also been found to encourage regrowth of cochlear hair cells in mice following genetic deletion or knock down with siRNA targeting p27. Research on hair cell regeneration may bring us closer to clinical treatment for human hearing loss caused by hair cell damage or death.
| Biology and health sciences | Sensory nervous system | Biology |
890259 | https://en.wikipedia.org/wiki/European%20mink | European mink | The European mink (Mustela lutreola), also known as the Russian mink and Eurasian mink, is a semiaquatic species of mustelid native to Europe.
It is similar in colour to the American mink, but is slightly smaller and has a less specialized skull. Despite having a similar name, build and behaviour, the European mink is not closely related to the American mink, being much closer to the European polecat and Siberian weasel (kolonok). The European mink occurs primarily by forest streams unlikely to freeze in winter. It primarily feeds on voles, frogs, fish, crustaceans and insects.
The European mink is listed by the IUCN as Critically Endangered due to an ongoing reduction in numbers, having been calculated as declining more than 50% over the past three generations and expected to decline at a rate exceeding 80% over the next three generations. European mink numbers began to shrink during the 19th century, with the species rapidly becoming extinct in some parts of Central Europe. During the 20th century, mink numbers declined all throughout their range, the reasons for which having been hypothesised to be due to a combination of factors, including climate change, competition with (as well as diseases spread by) the introduced American mink, habitat destruction, declines in crayfish numbers and hybridisation with the European polecat. In Central Europe and Finland, the decline preceded the introduction of the American mink, having likely been due to the destruction of river ecosystems, while in Estonia, the decline seems to coincide with the spread of the American mink.
Evolution and taxonomy
Fossil finds of the European mink are very rare, thus indicating the species is either a relative newcomer to Europe, probably having originated in North America, or a recent speciation caused by hybridization. It likely first arose in the Middle Pleistocene, with several fossils in Europe dated to the Late Pleistocene being found in caves and some suggesting early exploitation by humans. Genetic analyses indicate, rather than being closely related to the American mink, the European mink's closest relative is the European polecat (perhaps due to past hybridization) and the Siberian weasel, being intermediate in form between true polecats and other members of the genus. The closeness between the mink and polecat is emphasized by the fact the species can hybridize.
Subspecies
, seven subspecies are recognised.
Description
Build
The European mink is a typical representative of the genus Mustela, having a greatly elongated body with short limbs. However, compared to its close relative, the Siberian weasel, the mink is more compact and less thinly built, thus approaching ferrets and European polecats in build. The European mink has a large, broad head with short ears. The limbs are short, with relatively well-developed membranes between the digits, particularly on the hind feet. The mink's tail is short, and does not exceed half the animal's body length (constituting about 40% of its length). The European mink's skull is less elongated than the kolonok's, with more widely spaced zygomatic arches and has a less massive facial region. In general characteristics, the skull is intermediate in shape between that of the Siberian weasel and the European polecat. Overall, the skull is less specialized for carnivory than that of polecats and the American mink. Males measure in body length, while females measure . Tail length is in males and . Overall weight is . It is a fast and agile animal, which swims and dives skilfully. It is able to run along stream beds, and stay underwater for one to two minutes. When swimming, it paddles with both its front and back limbs simultaneously.
Fur
The winter fur of the European mink is very thick and dense, but not long, and quite loosely fitting. The underfur is particularly dense compared with that of more land-based members of the genus Mustela. The guard hairs are quite coarse and lustrous, with very wide contour hairs which are flat in the middle, as is typical in aquatic mammals. The length of the hairs on the back and belly differ little, a further adaptation to the European mink's semiaquatic way of life. The summer fur is somewhat shorter, coarser and less dense than the winter fur, though the differences are much less than in purely terrestrial mustelids.
In dark coloured individuals, the fur is dark brown or almost blackish-brown, while light individuals are reddish brown. Fur colour is evenly distributed over the whole body, though in a few cases, the belly is a bit lighter than the upper parts. In particularly dark individuals, a dark, broad dorsal belt is present. The limbs and tail are slightly darker than the trunk. The face has no colour pattern, though its upper and lower lips and chin are pure white. White markings may also occur on the lower surface of the neck and chest. Occasionally, colour mutations such as albinos and white spots throughout the pelage occur. The summer fur is somewhat lighter, and dirty in tone, with more reddish highlights.
Differences from American mink
The European mink is similar to the American mink, but with several important differences. The tail is longer in the American species, almost reaching half its body length. The winter fur of the American mink is denser, longer and more closely fitting than that of the European mink. Unlike the European mink, which has white patches on both upper and lower lips, the American mink almost never has white marks on the upper lip. The European mink's skull is much less specialised than the American species' in the direction of carnivory, bearing more infantile features, such as a weaker dentition and less strongly developed projections. The European mink is reportedly less efficient than the American species underwater.
Behaviour
Territorial and denning behaviours
The European mink does not form large territories, possibly due to the abundance of food on the banks of small water bodies. The size of each territory varies according to the availability of food; in areas with water meadows with little food, the home range is , though it is more usual for territories to be . Summer territories are smaller than winter territories. Along shorelines, the length of a home range varies from , with a width of .
The European mink has both a permanent burrow and temporary shelters. The former is used all year except during floods, and is located no more than from the water's edge. The construction of the burrow is not complex, often consisting of one or two passages in diameter and in length, leading to a nest chamber measuring . Nesting chambers are lined with straw, moss, mouse wool and bird feathers. It is more sedentary than the American mink, and will confine itself for long periods in its burrow in very cold weather.
Reproduction and development
During the mating season, the sexual organs of the female enlarge greatly and become pinkish-lilac in colour, which is in contrast with the female American mink, whose organs do not change. In the Moscow Zoo, estrus was observed on 22–26 April, with copulation lasting from 15 minutes to an hour. The average litter consists of three to seven kits. At birth, kits weigh about , and grow rapidly, trebling their weight 10 days after birth. They are born blind; the eyes open after 30–31 days. The lactation period lasts 2.0 to 2.5 months, though the kits eat solid food after 20–25 days. They accompany the mother on hunting expeditions at the age of 56–70 days, and become independent at the age 70–84 days.
Diet
The European mink has a diverse diet consisting largely of aquatic and riparian fauna. Differences between its diet and that of the American mink are small. Voles are the most important food source, closely followed by crustaceans, frogs and water insects. Fish are an important food source in floodlands, with cases being known of European minks catching fish weighing . The European mink's daily food requirement is . In times of food abundance, it caches its food.
Range and status
The European mink is mostly restricted to Europe. Its range was widespread in the 19th century, with a distribution extending from northern Spain in the west to the river Ob (just east of the Urals) in the east, and from the Archangelsk region in the north to the northern Caucasus in the south. Over the last 150 years, though, it has severely declined by more than 90% and been extirpated or greatly reduced over most of its former range. The current range includes an isolated population in northern Spain and western France, which is widely disjunct from the main range in Eastern Europe (Latvia, Estonia, Belarus, Ukraine, central regions of European Russia, the Danube Delta in Romania and northwestern Bulgaria). It occurs from sea level to . In Estonia, the European mink population has been successfully re-established on the island of Hiiumaa, and there are plans for repeating the process on the nearby island of Saaremaa.
Decline
The earliest actual records of decreases in European mink numbers occurred in Germany, having already become extinct in several areas by the middle of the 18th century. A similar pattern occurred in Switzerland, with no records of minks being published in the 20th century. Records of minks in Austria stopped by the late 18th century. By the 1930s–1950s, the European mink became extinct in Poland, Hungary, Czechoslovakia and possibly Bulgaria. In Finland, the main decline occurred in the 1920s-1950s and the species was thought to be extinct in the 1970s, though a few specimens were reported in the 1990s. In Latvia, the European mink was thought to be extinct for years, until a specimen was captured in 1992. In Lithuania, the last specimens were caught in 1978–79. The decline of the European mink in Estonia and Belarus was rapid during the 1980s, with only a few small, fragmented populations in the northeastern regions of both countries being reported in the 1990s. The decline of European mink numbers in Ukraine began in the late 1950s, with now only a few small and isolated populations being reported in the upper courses of the Ukrainian Carpathian rivers. Their numbers in Moldova began to drop very quickly in the 1930s, with the last known population having been confined to the lower course of the River Prut on the Romanian border by the late 1980s. In Romania, the European mink was very common and widely distributed, with 8000–10,000 being captured in 1960. Currently, Romanian mink populations are confined to the Danube Delta. In European Russia, the European mink was common and widespread in the early 20th century, but began to decline during the 1950s–1970s. The core of their range was in the Tver Region, though they began to decline there by the 1990s, which was worsened by a colonisation of the area by the American mink. Between 1981 and 1989, 388 European minks were introduced to two of the Kurile Islands, though by the 1990s, the population there was found to be lower than that originally released. In France and Spain, an isolated range occurs, extending from Brittany to northern Spain. Data from the 1990s indicate the European mink has disappeared from the northern half of this previous range.
Possible reasons for decline
Habitat loss
Habitat-related declines of European mink numbers may have started during the Little Ice Age, which was further aggravated by human activity. As the European mink is more dependent on wetland habitats than the American species, its decline in Central Europe, Estonia, Finland, Russia, Moldova and Ukraine has been linked to the drainage of small rivers. In mid-19th-century Germany, for example, European mink populations declined in a decade due to expanded land drainage. Although land improvement and river dredging certainly resulted in population decreases and fragmentation, in areas which still maintain suitable river ecosystems, such as Poland, Hungary, the former Czechoslovakia, Finland and Russia, the decline preceded the change in wetland habitats, and may have been caused by extensive agricultural development.
Overhunting
The European mink was historically hunted extensively, particularly in Russia, where in some districts, the decline prompted a temporary ban on mink hunting to let the population recover. In the early 20th century, 40–60,000 European minks were caught annually in the Soviet Union, with a record of 75,000 individuals (an estimate which exceeds the modern global European mink population). In Finland, annual mink catches reached 3000 specimens in the 1920s. In Romania, 10,000 minks were caught annually around 1960. However, this reason alone cannot account for the decline in areas where hunting was less intense, such as in Germany.
Decline of crayfish
The decline of European crayfish has been proposed as a factor in the drop in mink numbers, as minks are notably absent in the eastern side of the Urals, where crayfish are also absent. The decline in mink numbers has also been linked to the destruction of crayfish in Finland during the 1920s-1940s, when the crustaceans were infected with crayfish plague. The failure of the European mink to expand west to Scandinavia coincides with the gap in crayfish distribution.
Competition with the American mink and disease
The American mink was introduced and released in Europe during the 1920s–1930s. The American mink is less dependent on wetland habitats than the European mink and is 20-40% larger. The impact of feral American minks on European mink populations has been explained through the competitive exclusion principle and because the American mink reproduces a month earlier than the European species, and matings between male American minks and female European minks result in the embryos being reabsorbed. Thus, female European minks impregnated by male American minks are unable to reproduce with their conspecifics. Though the presence of the American mink has coincided with the decline of European mink numbers in Belarus and Estonia, the decline of the European mink in some areas preceded the introduction of the American mink by many years, and there are areas in Russia where the American species is absent, though European mink populations in these regions are still declining.
Diseases spread by the American mink can also account for the decline. Twenty-seven helminth species are recorded to infest the European mink, consisting of 14 trematodes, two cestodes and 11 nematodes. The mink is also vulnerable to pulmonary filariasis, krenzomatiasis and skrjabingylosis. In the Leningrad and Pskov Oblasts, 77.1% of European minks were found to be infected with skrjabingylosis.
Hybridisation and competition with the European polecat
In the early 20th century, northern Europe underwent a warm climatic period which coincided with an expansion of the range of the European polecat. The European mink possibly was gradually absorbed by the polecat due to hybridisation. Also, competition with the polecat has greatly increased, due to landscape change favouring the polecat. There is one record of a polecat attacking a mink and dragging it to its burrow.
Polecat-mink hybrids are termed khor'-tumak by furriers and khonorik by fanciers. Such hybridisation is very rare in the wild, and typically only occurs where European minks are declining. A polecat-mink hybrid has a poorly defined facial mask, yellow fur on the ears, grey-yellow underfur and long, dark brown guard hairs. Fairly large, the males attain the peak sizes known for European polecats (weighing and measuring in length), and females are much larger than female European minks (weighing and measuring in length). The majority of polecat-mink hybrids have skulls bearing greater similarities to those of polecats than to minks. Hybrids can swim well like minks and burrow for food like polecats. They are very difficult to tame and breed, as males are sterile, though females are fertile. The first captive polecat-mink hybrid was created in 1978 by Soviet zoologist Dr. Dmitry Ternovsky of Novosibirsk. Originally bred for their fur (which was more valuable than that of either parent species), the breeding of these hybrids declined as European mink populations decreased. Studies on the behavioural ecology of free-ranging polecat-mink hybrids in the upper reaches of the Lovat River indicate the hybrids will stray from aquatic habitats more readily than pure minks, and will tolerate both parent species entering their territories, though the hybrid's larger size (especially the male's) may deter intrusion. During the summer period, the diet of wild polecat-mink hybrids is more similar to that of the mink than to the polecat, as they feed predominantly on frogs. During the winter, their diets overlap more with those of polecats, and will eat a larger proportion of rodents than in the summer, though they still rely heavily on frogs and rarely scavenge ungulate carcasses as the polecat does.
Predation
Predators of the European mink include the European polecat, the American mink, the golden eagle, large owls and the red fox. Red fox numbers have increased greatly in areas where the wolf and Eurasian lynx have been extirpated, as well as areas where modern forestry is practised. As red foxes are known to prey on mustelids, excessive fox predation on the European mink is a possible factor, though it is improbable to have been a factor in Finland, where fox numbers were low during the early 20th century.
| Biology and health sciences | Mustelidae | Animals |
890314 | https://en.wikipedia.org/wiki/Brittle%20star | Brittle star | Brittle stars, serpent stars, or ophiuroids (; ; referring to the serpent-like arms of the brittle star) are echinoderms in the class Ophiuroidea, closely related to starfish. They crawl across the sea floor using their flexible arms for locomotion. The ophiuroids generally have five long, slender, whip-like arms which may reach up to in length on the largest specimens.
The Ophiuroidea contain two large clades, Ophiurida (brittle stars) and Euryalida (basket stars). Over 2,000 species of brittle stars live today. More than 1,200 of these species are found in deep waters, greater than 200 m deep.
Range
The ophiuroids diverged in the Early Ordovician. Ophiuroids can be found today in all of the major marine provinces, from the poles to the tropics. Basket stars are usually confined to the deeper parts of this range; Ophiuroids are known even from abyssal (>6,000 m) depths. However, brittle stars are also common members of reef communities, where they hide under rocks and even within other living organisms. A few ophiuroid species can even tolerate brackish water, an ability otherwise almost unknown among echinoderms. A brittle star's skeleton is made up of embedded ossicles.
Anatomy
Of all echinoderms, the Ophiuroidea may have the strongest tendency toward five-segment radial (pentaradial) symmetry. The body outline is similar to that of starfish, in that ophiuroids have five arms joined to a central body disk. However, in ophiuroids, the central body disk is sharply marked off from the arms.
The disk contains all of the viscera. That is, the internal organs of digestion and reproduction never enter the arms, as they do in the Asteroidea. The underside of the disk contains the mouth, which has five toothed jaws formed from skeletal plates. The madreporite is usually located within one of the jaw plates, and not on the upper side of the animal as it is in starfish.
The ophiuroid coelom is strongly reduced, particularly in comparison to other echinoderms.
Water-vascular system
The vessels of the water vascular system end in tube feet. The water vascular system generally has one madreporite. Others, such as certain Euryalina, have one per arm on the aboral surface. Still other forms have no madreporite at all. Suckers and ampullae are absent from the tube feet.
Nervous system
The nervous system consists of a main nerve ring which runs around the central disk. At the base of each arm, the ring attaches to a radial nerve which runs to the end of the limb. The nerves in each limb run through a canal at the base of the vertebral ossicles.
Most ophiuroids have no eyes, or other specialised sense organs. However, they have several types of sensitive nerve endings in their epidermis, and are able to sense chemicals in the water, touch, and even the presence or absence of light. Moreover, tube feet may sense light as well as odors. These are especially found at the ends of their arms, detecting light and retreating into crevices.
Digestion
The mouth is rimmed with five jaws, and serves as an anus (egestion) as well as a mouth (ingestion). Behind the jaws is a short esophagus and a stomach cavity which occupies much of the dorsal half of the disk. Digestion occurs within 10 pouches or infolds of the stomach, which are essentially ceca, but unlike in sea stars, almost never extend into the arms. The stomach wall contains glandular hepatic cells.
Ophiuroids are generally scavengers or detritivores. Small organic particles are moved into the mouth by the tube feet. Ophiuroids may also prey on small crustaceans or worms. Basket stars in particular may be capable of suspension feeding, using the mucus coating on their arms to trap plankton and bacteria. They extend one arm out and use the other four as anchors. Brittle stars will eat small suspended organisms if available. In large, crowded areas, brittle stars eat suspended matter from prevailing seafloor currents.
Many species in the family Ophiuridae are carnivorous. Ophiura Linnaeus hunts epibenthic animals and the Antarctic Ophiosparte gigas is an active predator. Ophiura albida Forbes and Ophiura sarsii Lütken eat both infaunal prey, carrion and seafloor organic matter, and Ophionereis reticulata is omnivorous and feeds on algae, polychaetes and detritus.
In basket stars, the arms are used to sweep food rhythmically to the mouth. Ophiopsammus maculata consumes Nothofagus pollen in the New Zealand fjords (since those trees hang over the water). Eurylina clings to coral branches to browse on the polyps.
Respiration
Gas exchange and excretion occur through cilia-lined sacs called bursae; each opens between the arm bases on the underside of the disk. Typically ten bursae are found, and each fits between two stomach digestive pouches. Water flows through the bursae by means of cilia or muscular contraction. Oxygen is transported through the body by the hemal system, a series of sinuses and vessels distinct from the water vascular system.
The bursae are probably also the main organs of excretion, with phagocytic "coelomocytes" collecting waste products in the body cavity and then migrating to the bursae for expulsion from the body.
Musculoskeletal system
Like all echinoderms, the Ophiuroidea possess a skeleton of calcium carbonate in the form of calcite. In ophiuroids, the calcite ossicles are fused to form armor plates which are known collectively as the test. The plates are covered by the epidermis, which consists of a smooth syncytium. In most species, the joints between the ossicles and superficial plates allow the arm to bend to the side, but cannot bend upwards. However, in the basket stars, the arms are flexible in all directions.
Both the Ophiurida and Euryalida (the basket stars) have five long, slender, flexible, whip-like arms, up to 60 cm in length. They are supported by an internal skeleton of calcium carbonate plates referred to as vertebral ossicles. These "vertebrae" articulate by means of ball-and-socket joints, and are controlled by muscles. They are essentially fused plates which correspond to the parallel ambulacral plates in sea stars and five Paleozoic families of ophiuroids. In modern forms, the vertebrae occur along the median of the arm.
The ossicles are surrounded by a relatively thin ring of soft tissue, and then by four series of jointed plates, one each on the upper, lower, and lateral surfaces of the arm. The two lateral plates often have a number of elongated spines projecting outwards; these help to provide traction against the substrate while the animal is moving. The spines, in ophiuroids, compose a rigid border to the arm edges, whereas in euryalids they are transformed into downward-facing clubs or hooklets. Euryalids are similar to ophiurids, if larger, but their arms are forked and branched. Ophiuroid podia generally function as sensory organs. They are not usually used for feeding, as in Asteroidea. In the Paleozoic era, brittle stars had open ambulacral grooves, but in modern forms, these are turned inward.
In living ophiuroids, the vertebrae are linked by well-structured longitudinal muscles. Ophiuroida moves horizontally, and Euryalina species moves vertically. The latter have bigger vertebrae and smaller muscles. They are less spasmodic, but can coil their arms around objects, holding on even after death. These movement patterns are distinct to the taxa, separating them. Ophiuroida moves quickly when disturbed. One arm presses ahead, whereas the other four act as two pairs of opposite levers, thrusting the body in a series of rapid jerks. Although adults do not use their tube feet for locomotion, very young stages use them as stilts and even serve as an adhesive structure.
Reproduction
The sexes are separate in most species, though a few are hermaphroditic or protandric. The gonads are located in the disk, and open into pouches between the arms, called genital bursae. Fertilization is external in most species, with the gametes being shed into the surrounding water through the bursal sacs. An exception is the Ophiocanopidae, in which the gonads do not open into bursae and are instead paired in a chain along the basal arm joints.
Many species brood developing larvae in the bursae, effectively giving birth to live young. A few, such as Amphipholus squamata, are truly viviparous, with the embryo receiving nourishment from the mother through the wall of the bursa. However, some species do not brood their young, and instead have a free-swimming larval stage. Referred to as an ophiopluteus, these larvae have four pairs of rigid arms lined with cilia. They develop directly into an adult, without the attachment stage found in most starfish larvae. The number of species exhibiting ophiopluteus larvae are fewer than those that directly develop.
In a few species, the female carries a dwarf male, clinging to it with the mouth.
Fission
Some brittle stars, such as the six-armed members of the family Ophiactidae, exhibit fissiparity (division through fission), with the disk splitting in half. Regrowth of both the lost part of the disk and the arms occur which yields an animal with three large arms and three small arms during the period of growth.
The West Indian brittle star, Ophiocomella ophiactoides, frequently undergoes asexual reproduction by fission of the disk with subsequent regeneration of the arms. In both summer and winter, large numbers of individuals with three long arms and three short arms can be found. Other individuals have half a disk and only three arms. A study of the age range of the population indicates little recruitment and fission is the primary means of reproduction in this species.
In this species, fission appears to start with the softening of one side of the disk and the initiation of a furrow. This deepens and widens until it extends across the disk and the animal splits in two. New arms begin to grow before the fission is complete, thus minimizing the time between possible successive divisions. The plane of fission varies so that some newly formed individuals have existing arms of different lengths. The time period between successive divisions is 89 days, so theoretically, each brittle star can produce 15 new individuals during the course of a year.
Life span
Brittle stars generally sexually mature in two to three years, become full grown in three to four years, and live up to five years. Members of Euryalina, such as Gorgonocephalus, may live much longer.
Regeneration
Ophiuroids can readily regenerate lost arms or arm segments unless all arms are lost. Ophiuroids use this ability to escape predators, in a way similar to lizards which deliberately shed the distal part of their tails to confuse pursuers. Moreover, the Amphiuridae can regenerate gut and gonad fragments lost along with the arms. Discarded arms have not been shown to have the ability to regenerate.
Locomotion
Brittle stars use their arms for locomotion. Brittle stars move fairly rapidly by wriggling their arms which are highly flexible and enable the animals to make either snake-like or rowing movements. However, they tend to attach themselves to the sea floor or to sponges or cnidarians, such as coral.
They move as if they were bilaterally symmetrical, with an arbitrary leg selected as the symmetry axis and the other four used in propulsion. The axial leg may be facing or trailing the direction of motion, and due to the radially symmetrical nervous system, can be changed whenever a change in direction is necessary.
Bioluminescence
Over 60 species of brittle stars are known to be bioluminescent. Most of these produce light in the green wavelengths, although a few blue-emitting species have also been discovered. Both shallow-water and deep-sea species of brittle stars are known to produce light. Presumably, this light is used to deter predators.
Ecology
Brittle stars live in areas from the low-tide level downwards. Six families live at least 2 m deep; the genera Ophiura, Amphiophiura, and Ophiacantha range below 4 m. Shallow species live among sponges, stones, or coral, or under the sand or mud, with only their arms protruding. Two of the best-known shallow species are the green brittle star (Ophioderma brevispina), found from Massachusetts to Brazil, and the common European brittle star (Ophiothrix fragilis). Deep-water species tend to live in or on the sea floor or adhere to coral, urchins, or xenophyophores. The most widespread species is the long-armed brittle star (Amphipholis squamata), a grayish or bluish, strongly luminescent species.
Parasites
The main parasite to enter the digestive tract or genitals are protozoans. Crustaceans, nematodes, trematodes, and polychaete annelids also serve as parasites. Algal parasites such as Coccomyxa ophiurae cause spinal malformation. Unlike in sea stars and sea urchins, annelids are not typical parasites.
Diversity and taxonomy
Between and species of brittle stars are currently known, but the total number of modern species may be over . This makes brittle stars the most abundant group of current echinoderms (before sea stars). Around 270 genera are known, these are distributed in 16 families, which makes them at the same time a relatively poorly diversified group structurally, compared with the other echinoderms. For example, 467 species belong to the sole family of Amphiuridae (frail brittle stars which live buried in the sediment leaving only their arms in the stream to capture the plankton). There are also 344 species in the family of Ophiuridae.
List of families according to the World Register of Marine Species, following O'Hara 2017:
subclass Myophiuroidea Matsumoto, 1915
infraclass Metophiurida Matsumoto, 1913
superorder Euryophiurida O'Hara, Hugall, Thuy, Stöhr & Martynov, 2017
order Euryalida Lamarck, 1816
family Asteronychidae Ljungman, 1867
family Euryalidae Gray, 1840
family Gorgonocephalidae Ljungman, 1867
order Ophiurida Müller & Troschel, 1840 sensu O'Hara et al., 2017
suborder Ophiomusina O'Hara et al., 2017
family Ophiomusaidae (O'Hara, Stöhr, Hugall, Thuy, Martynov, 2018)
family Ophiosphalmidae (O'Hara, Stöhr, Hugall, Thuy & Martynov, 2018)
Ophiomusina incertae sedis
suborder Ophiurina Müller & Troschel, 1840 sensu O'Hara et al., 2017
family Astrophiuridae Sladen, 1879
family Ophiopyrgidae Perrier, 1893
family Ophiuridae Müller & Troschel, 1840
Ophiurina incertae sedis
Ophiurida incertae sedis
superorder Ophintegrida O'Hara, Hugall, Thuy, Stöhr & Martynov, 2017
order Amphilepidida O'Hara, Hugall, Thuy, Stöhr & Martynov, 2017
suborder Gnathophiurina Matsumoto, 1915
superfamily Amphiuroidea Ljungman, 1867
family Amphiuridae Ljungman, 1867
family Amphilepididae Matsumoto, 1915
superfamily Ophiactoidea Ljungman, 1867
family Ophiactidae Matsumoto, 1915
family Ophiopholidae O'Hara, Stöhr, Hugall, Thuy & Martynov, 2018
family Ophiothamnidae O'Hara, Stöhr, Hugall, Thuy & Martynov, 2018
family Ophiotrichidae Ljungman, 1867
suborder Ophionereidina O'Hara, Hugall, Thuy, Stöhr & Martynov, 2017
superfamily Ophiolepidoidea Ljungman, 1867
family Hemieuryalidae Verrill, 1899
family Ophiolepididae Ljungman, 1867 (restricted)
superfamily Ophionereidoidea Ljungman, 1867
family Amphilimnidae O'Hara, Stöhr, Hugall, Thuy & Martynov, 2018
family Ophionereididae Ljungman, 1867
suborder Ophiopsilina Matsumoto, 1915
superfamily Ophiopsiloidea Matsumoto, 1915
family Ophiopsilidae Matsumoto, 1915
order Ophiacanthida O'Hara, Hugall, Thuy, Stöhr & Martynov, 2017
suborder Ophiacanthina O'Hara, Hugall, Thuy, Stöhr & Martynov, 2017
family Clarkcomidae O'Hara, Stöhr, Hugall, Thuy & Martynov, 2018
family Ophiacanthidae Ljungman, 1867
family Ophiobyrsidae Matsumoto, 1915
family Ophiocamacidae (O'Hara, Stöhr, Hugall, Thuy, Martynov, 2018)
family Ophiopteridae O'Hara, Stöhr, Hugall, Thuy & Martynov, 2018
family Ophiotomidae Paterson, 1985
family Ophiojuridae O'Hara, Thuy & Hugall, 2021
suborder Ophiodermatina Ljungman, 1867
superfamily Ophiocomoidea Ljungman, 1867
family Ophiocomidae Ljungman, 1867
superfamily Ophiodermatoidea Ljungman, 1867
family Ophiodermatidae Ljungman, 1867
family Ophiomyxidae Ljungman, 1867
family Ophiopezidae O'Hara, Stöhr, Hugall, Thuy & Martynov, 2018
Ophiacanthida incertae sedis
order Ophioleucida O'Hara, Hugall, Thuy, Stöhr & Martynov, 2017
family Ophiernidae O'Hara, Stöhr, Hugall, Thuy & Martynov, 2018
family Ophioleucidae Matsumoto, 1915
order Ophioscolecida O'Hara, Hugall, Thuy, Stöhr & Martynov, 2017
family Ophiohelidae Perrier, 1893
family Ophioscolecidae Lütken, 1869
Ophiuroidea incertae sedis
Fossil record
The first known brittle stars date from Early Ordovician. Study of past distribution and evolution of brittle stars has been hampered by the tendency of dead brittle stars to disarticulate and scatter, providing poor brittle star fossils. Until discoveries in the Agrio Formation of Neuquén Basin in the 2010s no fossil brittle star was known in the Southern Hemisphere, nor was any brittle star of Cretaceous age known.
Silurian fossils from a minor mass extinction called the Mulde event shows the ancestors of modern brittle stars went though a bottleneck, where a miniaturization caused by paedomorphosis led to structural simplification of their skeletal anatomy. These traits affected their further evolution. As they began to increase in size again, so did their complexity. The first large-sized modern brittle star originated in the Early Carboniferous.
Human relations
Brittle stars are not used as food, though they are not toxic, because of their strong skeleton.
Even if some species have blunt spines, no brittlestar is known to be dangerous, nor venomous. There is no harm evidence towards humans, and even with their predators, brittlestars' only means of defense is escaping or discarding an arm.
Aquaria
Brittle stars are a moderately popular invertebrate in fishkeeping. They can easily thrive in marine tanks; in fact, the micro brittle star is a common "hitchhiker" that will propagate and become common in almost any saltwater tank, if one happens to come along on some live rock.
Larger brittle stars are popular because, unlike Asteroidea, they are not generally seen as a threat to coral, and are also faster-moving and more active than their more archetypical cousins.
| Biology and health sciences | Echinoderms | Animals |
20159769 | https://en.wikipedia.org/wiki/Droop%20speed%20control | Droop speed control | Droop speed control is a control mode used for AC electrical power generators, whereby the power output of a generator reduces as the line frequency increases. It is commonly used as the speed control mode of the governor of a prime mover driving a synchronous generator connected to an electrical grid. It works by controlling the rate of power produced by the prime mover according to the grid frequency. With droop speed control, when the grid is operating at maximum operating frequency, the prime mover's power is reduced to zero, and when the grid is at minimum operating frequency, the power is set to 100%, and intermediate values at other operating frequencies.
This mode allows synchronous generators to run in parallel, so that loads are shared among generators with the same droop curve in proportion to their power rating.
In practice, the droop curves that are used by generators on large electrical grids are not necessarily linear or the same, and may be adjusted by operators. This permits the ratio of power used to vary depending on load, so for example, base load generators will generate a larger proportion at low demand. Stability requires that over the operating frequency range the power output is a monotonically decreasing function of frequency.
Droop speed control can also be used by grid storage systems. With droop speed control those systems will remove energy from the grid at higher than average frequencies, and supply it at lower frequencies.
Linear
The frequency of a synchronous generator is given by
where
F, frequency (in Hz),
P, number of poles,
N, speed of generator (in RPM)
The frequency (F) of a synchronous generator is directly proportional to its speed (N). When multiple synchronous generators are connected in parallel to the electrical grid, the frequency is fixed by the grid, since individual power output of each generator will be small compared to the load on a large grid. Synchronous generators connected to the grid run at various speeds but they all run at the same frequency because they differ in the number of poles (P).
A speed reference as percentage of actual speed is set in this mode. As the generator is loaded from no load to full load, the actual speed of the prime mover tends to decrease. In order to increase the power output in this mode, the prime mover speed reference is increased. Because the actual prime mover speed is fixed by the grid, this difference in speed reference and actual speed of the prime mover is used to increase the flow of working fluid (fuel, steam, etc.) to the prime mover, and hence power output is increased. The reverse will be true for decreasing power output. The prime mover speed reference is always greater than actual speed of the prime mover. The actual speed of the prime mover is allowed to "droop" or decrease with respect to the reference, and so the name.
For example, if the turbine is rated at 3000 rpm, and the machine speed reduces from 3000 rpm to 2880 rpm when it is loaded from no load to base load, then the droop % is given by
= (3000 – 2880) / 3000
= 4%
In this case, speed reference will be 104% and actual speed will be 100%. For every 1% change in the turbine speed reference, the power output of the turbine will change by 25% of rated for a unit with a 4% droop setting. Droop is therefore expressed as the percentage change in (design) speed required for 100% governor action.
As frequency is fixed on the grid, and so actual turbine speed is also fixed, the increase in turbine speed reference will increase the error between reference and actual speed. As the difference increases, fuel flow is increased to increase power output, and vice versa. This type of control is referred to as "straight proportional" control. If the entire grid tends to be overloaded, the grid frequency and hence actual speed of generator will decrease. All units will see an increase in the speed error, and so increase fuel flow to their prime movers and power output. In this way droop speed control mode also helps to hold a stable grid frequency. The amount of power produced is strictly proportional to the error between the actual turbine speed and speed reference.
It can be mathematically shown that if all machines synchronized to a system have the same droop speed control, they will share load proportionate to the machine ratings.
For example, how fuel flow is increased or decreased in a GE-design heavy duty gas turbine can be given by the formula,
FSRN = (FSKRN2 * (TNR-TNH)) + FSKRN1
Where,
FSRN = Fuel Stroke Reference (Fuel supplied to Gas Turbine) for droop mode
TNR = Turbine Speed Reference
TNH = Actual Turbine Speed
FSKRN2 = Constant
FSKRN1 = Constant
The above formula is nothing but the equation of a straight line (y = mx + b).
Multiple synchronous generators having equal % droop setting connected to a grid will share the change in grid load in proportion of their base load.
For stable operation of the electrical grid of North America, power plants typically operate with a four or five percent speed droop. By definition, with 5% droop the full-load speed is 100% and the no-load speed is 105%.
Normally the changes in speed are minor due to inertia of the total rotating mass of all generators and motors running on the grid. Adjustments in power output for a particular primer mover and generator combination are made by slowly raising the droop curve by increasing the spring pressure on a centrifugal governor or by an engine control unit adjustment, or the analogous operation for an electronic speed governor. All units to be connected to a grid should have the same droop setting, so that all plants respond in the same way to the instantaneous changes in frequency without depending on outside communication.
Next to the inertia given by the parallel operation of synchronous generators, the frequency speed droop is the primary instantaneous parameter in control of an individual power plant's power output (kW).
| Technology | Concepts | null |
1996352 | https://en.wikipedia.org/wiki/Bridge%20of%20the%20Americas | Bridge of the Americas | The Bridge of the Americas (; originally known as the Thatcher Ferry Bridge) is a road bridge in Panama which spans the Pacific entrance to the Panama Canal. Designed by Sverdrup & Parcel, it was completed in 1962 at a cost of US$20 million, connecting the north and south American land masses (hence its name), connecting the American Continent. Two other bridges cross the canal: the Atlantic Bridge at the Gatun locks and the Centennial Bridge.
Description
The Bridge of the Americas crosses the Pacific approach to the Panama Canal at Balboa, near Panama City. It was built between 1959 and 1962 by the United States at a cost of US$20 million. From its completion in 1962 until the opening of the parallel Centennial Bridge in 2004, the Bridge of the Americas was a key part of the Pan-American Highway. The Bridge of the Americas greatly increases road traffic capacity across the Canal. Two earlier bridges cross the Canal, but they use moveable designs and have limited traffic capacity. These earlier spans include a small swinging road bridge, built into the lock structure at Gatún, and a swinging road/rail bridge constructed in 1942 at Miraflores. The Centennial Bridge was constructed in an effort to eliminate the bottleneck of, and reduce traffic congestion on, the Bridge of the Americas.
The bridge is a cantilever design where the suspended span is a tied arch. The bridge has a total length of 1,654 m (5,425 ft) in 14 spans, abutment to abutment. The main span measures and the tied arch (the center part of the main span) is . The highest point of the bridge is above mean sea level; the clearance under the main span is at high tide. Ships must cross under this bridge when traversing the Panama Canal, and are subject to this height restriction. The world's largest cruise ships, Oasis of the Seas, Allure of the Seas, Harmony of the Seas and the Symphony of the Seas will fit within the canal's widened locks, but they are too tall to pass under the Bridge of the Americas, even at low tide, unless the Bridge of the Americas is either raised or replaced in the future. (The Centennial Bridge is also a fixed obstacle, but its clearance is much higher: .)
The bridge has wide access ramps at each end, and pedestrian walkways on each side.
History
The need for a bridge
From the beginning of the French project to construct a canal, it was recognised that the cities of Colón and Panamá would be split from the rest of the republic by the new canal. This was an issue even during construction, when barges were used to ferry construction workers across the canal.
After the canal opened, the increasing number of cars, and the construction of a new road leading to Chiriquí, in the west of Panama, increased the need for some kind of crossing. The Panama Canal Mechanical Division addressed this in August 1931, with the commissioning of two new ferries, the Presidente Amador and President Washington. This service was expanded in August 1940, with additional barges mainly serving the military.
On June 3, 1942, a road/rail swing bridge was inaugurated at the Miraflores locks; although only usable when no ships were passing, this provided some relief for traffic wishing to cross the canal. Still, it was clear that a more substantial solution would be required. To meet the growing needs of vehicle traffic, another ferry, the Presidente Porras, was added in November 1942.
The bridge project
The idea of a permanent bridge over the canal had been proposed as a major priority as early as 1923. Subsequent administrations of Panama pressed this issue with the United States, which controlled the Canal Zone; and in 1955 the Remón-Eisenhower treaty committed the United States to building a bridge.
A contract worth $20,000,000 was awarded to John F. Beasly & Company who built the bridge out of steel and reinforced concrete, and the project was initiated in a ceremony which took place on December 23, 1958, in the presence of United States Ambassador Julian Harrington, and Panamanian President Ernesto de la Guardia Navarro. Construction began on October 12, 1959, and took nearly two and a half years to complete.
The inauguration of the bridge took place on October 12, 1962, with great ceremony. The ribbon was cut by Maurice H. Thatcher, after which those present were allowed to walk across the bridge. The ceremony was given full nationwide coverage on radio and television; significant precautions were taken to manage the large crowds of people present. These proved inadequate, however, and pro-Panamanian protesters disrupted the ceremony, even removing the memorial plaques on the bridge.
Post-construction
When opened, the bridge was an important part of the Pan-American Highway, and carried around 9,500 vehicles per day; however, this expanded over time, and by 2004 the bridge was carrying 35,000 vehicles per day. The bridge therefore became a significant bottleneck on the highway, which led to the construction of the Centennial Bridge, which now carries the Pan-American Highway too. On May 18, 2010, the bulk cargo ship Atlantic Hero struck one of the protective bases of the bridge after losing engine power, partially blocking that section of the canal to shipping traffic. The bridge did not receive damage and there were no fatalities. In December 2010, the Centennial Bridge access road collapsed in a mudslide, and commercial traffic was diverted to the Bridge of The Americas.
The "Thatcher Ferry Bridge"
The bridge was originally named "Thatcher Ferry Bridge", after the original ferry which crossed the canal at about the same point. The ferry was, in turn, named after Maurice H. Thatcher, a former member of the Canal Commission, who introduced the legislation which created the ferry. Thatcher cut the tape at the inauguration of the bridge.
The name was unpopular with the government of Panama, however, which preferred the name "Bridge of the Americas". The Panamanian view was made official by a resolution of the National Assembly on October 2, 1962, ten days before the inauguration.
The resolution read as follows:
The bridge over the Panama Canal shall bear the name Bridge of the Americas. Said name will be used exclusively to identify said bridge.
Panamanian government officials shall reject any document in which reference is made to the bridge by any name other than "Bridge of the Americas".
A copy of this resolution, with the appropriate note on style, shall be forwarded to all legislative bodies of the world, so that all may give the bridge the name chosen by this honorable assembly, complying with the express will of the Panamanian people.
Given in the city of Panama on the second day of the month of October of nineteen hundred and sixty-two.
President, Jorge Rubén Rosas
Secretary, Alberto Arango N.
During the inauguration ceremony (which was concluded with the playing of the "Thatcher Ferry Bridge March"), U.S. Under Secretary of State George Wildman Ball said in his speech: "we can look today to this bridge as a new and bright step toward the realization of that dream of a Pan-American Highway, which is now almost a reality. The grand bridge we inaugurate today — truly a bridge of the Americas — completes the last stage of the highway from the United States to Panama".
Nonetheless, the official name of the bridge became the "Thatcher Ferry Bridge" and remained so until Panamanian control in 1979.
Postage stamps were issued with the name "Thatcher Ferry Bridge". In the postage stamps and postal history of the Canal Zone they are well known for an error on one sheet where the bridge is missing.
| Technology | Bridges | null |
1996536 | https://en.wikipedia.org/wiki/Grain%20boundary | Grain boundary | In materials science, a grain boundary is the interface between two grains, or crystallites, in a polycrystalline material. Grain boundaries are two-dimensional defects in the crystal structure, and tend to decrease the electrical and thermal conductivity of the material. Most grain boundaries are preferred sites for the onset of corrosion and for the precipitation of new phases from the solid. They are also important to many of the mechanisms of creep. On the other hand, grain boundaries disrupt the motion of dislocations through a material, so reducing crystallite size is a common way to improve mechanical strength, as described by the Hall–Petch relationship.
High and low angle boundaries
It is convenient to categorize grain boundaries according to the extent of misorientation between the two grains. Low-angle grain boundaries (LAGB) or subgrain boundaries are those with a misorientation less than about 15 degrees. Generally speaking they are composed of an array of dislocations and their properties and structure are a function of the misorientation. In contrast the properties of high-angle grain boundaries, whose misorientation is greater than about 15 degrees (the transition angle varies from 10 to 15 degrees depending on the material), are normally found to be independent of the misorientation. However, there are 'special boundaries' at particular orientations whose interfacial energies are markedly lower than those of general high-angle grain boundaries.
The simplest boundary is that of a tilt boundary where the rotation axis is parallel to the boundary plane. This boundary can be conceived as forming from a single, contiguous crystallite or grain which is gradually bent by some external force. The energy associated with the elastic bending of the lattice can be reduced by inserting a dislocation, which is essentially a half-plane of atoms that act like a wedge, that creates a permanent misorientation between the two sides. As the grain is bent further, more and more dislocations must be introduced to accommodate the deformation resulting in a growing wall of dislocations – a low-angle boundary. The grain can now be considered to have split into two sub-grains of related crystallography but notably different orientations.
An alternative is a twist boundary where the misorientation occurs around an axis that is perpendicular to the boundary plane. This type of boundary incorporates two sets of screw dislocations. If the Burgers vectors of the dislocations are orthogonal, then the dislocations do not strongly interact and form a square network. In other cases, the dislocations may interact to form a more complex hexagonal structure.
These concepts of tilt and twist boundaries represent somewhat idealized cases. The majority of boundaries are of a mixed type, containing dislocations of different types and Burgers vectors, in order to create the best fit between the neighboring grains.
If the dislocations in the boundary remain isolated and distinct, the boundary can be considered to be low-angle. If deformation continues, the density of dislocations will increase and so reduce the spacing between neighboring dislocations. Eventually, the cores of the dislocations will begin to overlap and the ordered nature of the boundary will begin to break down. At this point the boundary can be considered to be high-angle and the original grain to have separated into two entirely separate grains.
In comparison to low-angle grain boundaries, high-angle boundaries are considerably more disordered, with large areas of poor fit and a comparatively open structure. Indeed, they were originally thought to be some form of amorphous or even liquid layer between the grains. However, this model could not explain the observed strength of grain boundaries and, after the invention of electron microscopy, direct evidence of the grain structure meant the hypothesis had to be discarded. It is now accepted that a boundary consists of structural units which depend on both the misorientation of the two grains and the plane of the interface. The types of structural unit that exist can be related to the concept of the coincidence site lattice, in which repeated units are formed from points where the two misoriented \
In coincident site lattice (CSL) theory, the degree of fit (Σ) between the structures of the two grains is described by the reciprocal of the ratio of coincidence sites to the total number of sites.
In this framework, it is possible to draw the lattice for the two grains and count the number of atoms that are shared (coincidence sites), and the total number of atoms on the boundary (total number of site). For example, when Σ=3 there will be one atom of each three that will be shared between the two lattices. Thus a boundary with high Σ might be expected to have a higher energy than one with low Σ. Low-angle boundaries, where the distortion is entirely accommodated by dislocations, are Σ1. Some other low-Σ boundaries have special properties, especially when the boundary plane is one that contains a high density of coincident sites. Examples include coherent twin boundaries (e.g., Σ3) and high-mobility boundaries in FCC materials (e.g., Σ7). Deviations from the ideal CSL orientation may be accommodated by local atomic relaxation or the inclusion of dislocations at the boundary.
Describing a boundary
A boundary can be described by the orientation of the boundary to the two grains and the 3-D rotation required to bring the grains into coincidence. Thus a boundary has 5 macroscopic degrees of freedom. However, it is common to describe a boundary only as the orientation relationship of the neighbouring grains. Generally, the convenience of ignoring the boundary plane orientation, which is very difficult to determine, outweighs the reduced information.
The relative orientation of the two grains is described using the rotation matrix:
Using this system the rotation angle θ is:
while the direction [uvw] of the rotation axis is:
The nature of the crystallography involved limits the misorientation of the boundary. A completely random polycrystal, with no texture, thus has a characteristic distribution of boundary misorientations (see figure). However, such cases are rare and most materials will deviate from this ideal to a greater or lesser degree.
Boundary energy
The energy of a low-angle boundary is dependent on the degree of misorientation between the neighbouring grains up to the transition to high-angle status. In the case of simple tilt boundaries the energy of a boundary made up of dislocations with Burgers vector b and spacing h is predicted by the Read–Shockley equation:
where:
with is the shear modulus, is Poisson's ratio, and is the radius of the dislocation core. It can be seen that as the energy of the boundary increases the energy per dislocation decreases. Thus there is a driving force to produce fewer, more misoriented boundaries (i.e., grain growth).
The situation in high-angle boundaries is more complex. Although theory predicts that the energy will be a minimum for ideal CSL configurations, with deviations requiring dislocations and other energetic features, empirical measurements suggest the relationship is more complicated. Some predicted troughs in energy are found as expected while others missing or substantially reduced.
Surveys of the available experimental data have indicated that simple relationships such as low are misleading:
It is concluded that no general and useful criterion for low energy can be enshrined in a simple geometric framework. Any understanding of the variations of interfacial energy must take account of the atomic structure and the details of the bonding at the interface.
Excess volume
The excess volume is another important property in the characterization of grain boundaries. Excess volume was first proposed by Bishop in a private communication to Aaron and Bolling in 1972. It describes how much expansion is induced by the presence of a GB and is thought that the degree and susceptibility of segregation is directly proportional to this. Despite the name the excess volume is actually a change in length, this is because of the 2D nature of GBs the length of interest is the expansion normal to the GB plane. The excess volume () is defined in the following way,
at constant temperature , pressure and number of atoms . Although a rough linear relationship between GB energy and excess volume exists the orientations where this relationship is violated can behave significantly differently affecting mechanical and electrical properties.
Experimental techniques have been developed which directly probe the excess volume and have been used to explore the properties of nanocrystalline copper and nickel. Theoretical methods have also been developed and are in good agreement. A key observation is that there is an inverse relationship with the bulk modulus meaning that the larger the bulk modulus (the ability to compress a material) the smaller the excess volume will be, there is also direct relationship with the lattice constant, this provides methodology to find materials with a desirable excess volume for a specific application.
Boundary migration
The movement of grain boundaries (HAGB) has implications for recrystallization and grain growth while subgrain boundary (LAGB) movement strongly influences recovery and the nucleation of recrystallization.
A boundary moves due to a pressure acting on it. It is generally assumed that the velocity is directly proportional to the pressure with the constant of proportionality being the mobility of the boundary. The mobility is strongly temperature dependent and often follows an Arrhenius type relationship:
The apparent activation energy (Q) may be related to the thermally activated atomistic processes that occur during boundary movement. However, there are several proposed mechanisms where the mobility will depend on the driving pressure and the assumed proportionality may break down.
It is generally accepted that the mobility of low-angle boundaries is much lower than that of high-angle boundaries. The following observations appear to hold true over a range of conditions:
The mobility of low-angle boundaries is proportional to the pressure acting on it.
The rate controlling process is that of bulk diffusion
The boundary mobility increases with misorientation.
Since low-angle boundaries are composed of arrays of dislocations and their movement may be related to dislocation theory. The most likely mechanism, given the experimental data, is that of dislocation climb, rate limited by the diffusion of solute in the bulk.
The movement of high-angle boundaries occurs by the transfer of atoms between the neighbouring grains. The ease with which this can occur will depend on the structure of the boundary, itself dependent on the crystallography of the grains involved, impurity atoms and the temperature. It is possible that some form of diffusionless mechanism (akin to diffusionless phase transformations such as martensite) may operate in certain conditions. Some defects in the boundary, such as steps and ledges, may also offer alternative mechanisms for atomic transfer.
Since a high-angle boundary is imperfectly packed compared to the normal lattice it has some amount of free space or free volume where solute atoms may possess a lower energy. As a result, a boundary may be associated with a solute atmosphere that will retard its movement. Only at higher velocities will the boundary be able to break free of its atmosphere and resume normal motion.
Both low- and high-angle boundaries are retarded by the presence of particles via the so-called Zener pinning effect. This effect is often exploited in commercial alloys to minimise or prevent recrystallization or grain growth during heat-treatment.
Complexion
Grain boundaries are the preferential site for segregation of impurities, which may form a thin layer with a different composition from the bulk and a variety of atomic structures that are distinct from the abutting crystalline phases. For example, a thin layer of silica, which also contains impurity cations, is often present in silicon nitride.
Grain boundary complexions were introduced by Ming Tang, Rowland Cannon, and W. Craig Carter in 2006.
These grain boundary phases are thermodynamically stable and can be considered as quasi-two-dimensional phase, which may undergo to transition, similar to those of bulk phases. In this case structure and chemistry abrupt changes are possible at a critical value of a thermodynamic parameter like temperature or pressure. This may strongly affect the macroscopic properties of the material, for example the electrical resistance or creep rates. Grain boundaries can be analyzed using equilibrium thermodynamics but cannot be considered as phases, because they do not satisfy Gibbs' definition: they are inhomogeneous, may have a gradient of structure, composition or properties. For this reasons they are defined as complexion: an interfacial material or stata that is in thermodynamic equilibrium with its abutting phases, with a finite and stable thickness (that is typically 2–20 Å). A complexion need the abutting phase to exist and its composition and structure need to be different from the abutting phase. Contrary to bulk phases, complexions also depend on the abutting phase. For example, silica rich amorphous layer present in Si3N3, is about 10 Å thick, but for special boundaries this equilibrium thickness is zero. Complexion can be grouped in 6 categories, according to their thickness: monolayer, bilayer, trilayer, nanolayer (with equilibrium thickness between 1 and 2 nm) and wetting. In the first cases the thickness of the layer will be constant; if extra material is present it will segregate at multiple grain junction, while in the last case there is no equilibrium thickness and this is determined by the amount of secondary phase present in the material. One example of grain boundary complexion transition is the passage from dry boundary to biltilayer in Au-doped Si, which is produced by the increase of Au.
Effect to the electronic structure
Grain boundaries can cause failure mechanically by embrittlement through solute segregation (see Hinkley Point A nuclear power station) but they also can detrimentally affect the electronic properties. In metal oxides it has been shown theoretically that at the grain boundaries in Al2O3 and MgO the insulating properties can be significantly diminished. Using density functional theory computer simulations of grain boundaries have shown that the band gap can be reduced by up to 45%. In the case of metals grain boundaries increase the resistivity as the size of the grains relative to the mean free path of other scatters becomes significant.
Defect concentration near grain boundaries
It is known that most materials are polycrystalline and contain grain boundaries and that grain boundaries can act as sinks and transport pathways for point defects. However experimentally and theoretically determining what effect point defects have on a system is difficult. Interesting examples of the complications of how point defects behave has been manifested in the temperature dependence of the Seebeck effect. In addition the dielectric and piezoelectric response can be altered by the distribution of point defects near grain boundaries. Mechanical properties can also be significantly influenced with properties such as the bulk modulus and damping being influenced by changes to the distribution of point defects within a material. It has also been found that the Kondo effect within graphene can be tuned due to a complex relationship between grain boundaries and point defects. Recent theoretical calculations have revealed that point defects can be extremely favourable near certain grain boundary types and significantly affect the electronic properties with a reduction in the band gap.
Relationship between theory and experiment
There has been a significant amount of work experimentally to observe both the structure and measure the properties of grain boundaries but the five dimensional degrees of freedom of grain boundaries within complex polycrystalline networks has not yet been fully understood and thus there is currently no method to control the structure and properties of most metals and alloys with atomic precision. Part of the problem is related to the fact that much of the theoretical work to understand grain boundaries is based upon construction of bicrystal (two) grains which do not represent the network of grains typically found in a real system and the use of classical force fields such as the embedded atom method often do not describe the physics near the grains correctly and density functional theory could be required to give realistic insights. Accurate modelling of grain boundaries both in terms of structure and atomic interactions could have the effect of improving engineering which could reduce waste and increase efficiency in terms of material usage and performance. From a computational point of view much of the research on grain boundaries has focused on bi-crystal systems, these are systems which only consider two grain boundaries. There has been recent work which has made use of novel grain evolution models which show that there are substantial differences in the material properties associated with whether curved or planar grains are present.
| Physical sciences | Crystallography | Physics |
1996903 | https://en.wikipedia.org/wiki/Cassiopeia%20A | Cassiopeia A | Cassiopeia A (Cas A) () is a supernova remnant (SNR) in the constellation Cassiopeia and the brightest extrasolar radio source in the sky at frequencies above 1 GHz. The supernova occurred approximately away within the Milky Way; given the width of the Orion Arm, it lies in the next-nearest arm outwards, the Perseus Arm, about 30 degrees from the Galactic anticenter. The expanding cloud of material left over from the supernova now appears approximately across from Earth's perspective. It has been seen in wavelengths of visible light with amateur telescopes down to 234 mm (9.25 in) with filters.
It is estimated that light from the supernova itself first reached Earth near the 1660s, although there are no definitively corresponding records from then. Cas A is circumpolar at and above mid-Northern latitudes which had extensive records and basic telescopes. Its likely omission in records is probably due to interstellar dust absorbing optical wavelength radiation before it reached Earth, although it is possible that it was recorded as a sixth magnitude star 3 Cassiopeiae by John Flamsteed. Possible explanations lean toward the idea that the source star was unusually massive and had previously ejected much of its outer layers. These outer layers would have cloaked the star and absorbed much of the visible-light emission as the inner star collapsed.
Cas A was among the first discrete astronomical radio sources found. Its discovery was reported in 1948 by Martin Ryle and Francis Graham-Smith, astronomers at Cambridge, based on observations with the Long Michelson Interferometer. The optical component was first identified in 1950.
Possible observations
Calculations working back from the currently observed expansion point to an explosion that would have become visible on Earth around 1667. Astronomer William Ashworth and others have suggested that the Astronomer Royal John Flamsteed may have inadvertently observed the supernova on , when he catalogued a sixth-magnitude star 3 Cassiopeiae, but there is no corresponding star at the recorded position. Possible explanations include an error in the position, or that a transient was recorded. Caroline Herschel noted that a star in the vicinity of τ Cas, HD 220562, fit well with 3 Cas if a common error in sextant readings was made. Alternatively, the star AR Cassiopeiae may have been observed, again with the position recorded incorrectly. The position and timing mean that it may have been an observation of the Cassiopeia A progenitor supernova. Another suggestion from recent cross-disciplinary research is that the supernova was the "noon day star", observed in 1630, that was thought to have heralded the birth of Charles II, the future monarch of Great Britain.
However, it is more probable that the "noon day star" was the planet Venus that reached its maximum morning brightness two days earlier, allowing day time visibility in a clear sky. A bright supernova in Cassiopeia would have been visible for months and there would be more observation records as Cassiopeia is visible above the horizon any night in Europe.
No supernova occurring within the Milky Way has been visible to the naked eye from Earth since.
Expansion
The expansion shell has a temperature of around 30 million K, and is expanding at 4,000−6,000 km/s.
Observations of the exploded star through the Hubble Space Telescope have shown that, despite the original belief that the remnants were expanding in a uniform manner, there are high velocity outlying eject knots moving with transverse velocities of 5,500−14,500 km/s with the highest speeds occurring in two nearly opposing jets. When the view of the expanding star uses colors to differentiate materials of different chemical compositions, it shows that similar materials often remain gathered together in the remnants of the explosion.
Radio source
Cas A had a flux density of at 1 GHz in 1980. Because the supernova remnant is cooling, its flux density is decreasing. At 1 GHz, its flux density is decreasing at a rate of per year. This decrease means that, at frequencies below 1 GHz, Cas A is now less intense than Cygnus A. Cas A is still the brightest extrasolar radio source in the sky at frequencies above 1 GHz.
X-ray source
Although Cas X-1 (or Cas XR-1), the apparent first X-ray source in the constellation Cassiopeia was not detected during the 16 June 1964, Aerobee sounding rocket flight, it was considered as a possible source. Cas A was scanned during another Aerobee rocket flight of 1 October 1964, but no significant X-ray flux above background was associated with the position. Cas XR-1 was discovered by an Aerobee rocket flight on 25 April 1965, at RA Dec . Cas X-1 is Cas A, a Type II SNR at RA Dec .
The designations Cassiopeia X-1, Cas XR-1, Cas X-1 are no longer used, but the X-ray source is Cas A (SNR G111.7-02.1) at 2U 2321+58.
In 1999, the Chandra X-Ray Observatory found CXOU J232327.8+584842, a central compact object that is the neutron star remnant left by the explosion.
Supernova reflected echo
In 2005 an infrared echo of the Cassiopeia A explosion was observed on nearby gas clouds using Spitzer Space Telescope. The infrared echo was also seen by IRAS and studied with the Infrared Spectrograph. Previously it was suspected that a flare in 1950 from a central pulsar could be responsible for the infrared echo. With the new data it was concluded that this is unlikely the case and that the infrared echo was caused by thermal emission by dust, which was heated by the radiative output of the supernova during the shock breakout. The infrared echo is accompanied by a scattered light echo. The recorded spectrum of the optical light echo proved the supernova was of Type IIb, meaning it resulted from the internal collapse and violent explosion of a massive star, most probably a red supergiant with a helium core which had lost almost all of its hydrogen envelope. This was the first observation of the light echo of a supernova whose explosion had not been directly observed which opens up the possibility of studying and reconstructing past astronomical events. In 2011 a study used spectra from different positions of the light echo to confirm that the Cassiopeia A supernova was asymmetric.
Phosphorus detection
In 2013, astronomers detected phosphorus in Cassiopeia A, which confirmed that this element is produced in supernovae through supernova nucleosynthesis. The phosphorus-to-iron ratio in material from the supernova remnant could be up to 100 times higher than in the Milky Way in general.
Gallery
| Physical sciences | Notable nebulae | Astronomy |
1997989 | https://en.wikipedia.org/wiki/Crystal%20twinning | Crystal twinning | Crystal twinning occurs when two or more adjacent crystals of the same mineral are oriented so that they share some of the same crystal lattice points in a symmetrical manner. The result is an intergrowth of two separate crystals that are tightly bonded to each other. The surface along which the lattice points are shared in twinned crystals is called a composition surface or twin plane.
Crystallographers classify twinned crystals by a number of twin laws, which are specific to the crystal structure. The type of twinning can be a diagnostic tool in mineral identification. There are three main types of twinning. The first is growth twinning which can occur both in very large and very small particles. The second is transformation twinning, where there is a change in the crystal structure. The third is deformation twinning, in which twinning develops in a crystal in response to a shear stress, and is an important mechanism for permanent shape changes in a crystal.
Definition
Twinning, a version of macle, is a form of symmetrical intergrowth between two or more adjacent crystals of the same mineral. It differs from the ordinary random intergrowth of mineral grains in a mineral deposit, because the relative orientations of the two crystal segments show a fixed relationship that is characteristic of the mineral structure. The relationship is defined by a symmetry operation called a twin operation.
The twin operation is not one of the normal symmetry operations of the untwinned crystal structure. For example, the twin operation may be reflection across a plane that is not a symmetry plane of the single crystal.
On the microscopic level, the twin boundary is characterized by a set of atomic positions in the crystal lattice that are shared between the two orientations. These shared lattice points give the junction between the crystal segments much greater strength than that between randomly oriented grains, so that the twinned crystals do not easily break apart.
Parallel growth describes a form of crystal growth that produces the appearance of a cluster of aligned crystals which could be mistaken for twins. Close examination reveals that the cluster is actually a single crystal. This is not twinning, since the crystal lattice is continuous throughout the cluster. Parallel growth likely takes place because it reduces system energy.
Twin laws
Twin laws are symmetry operations that define the orientation between twin crystal segments. These are as characteristic of the mineral as are its crystal face angles. For example, crystals of staurolite show twinning at angles of almost precisely 90 degrees or 30 degrees. A twin law is not a symmetry operation of the full set of basis points.
Twin laws include reflection operations, rotation operations, and the inversion operation. Reflection twinning is described by the Miller indices of the twin plane (i.e. {hkl}) while rotational twinning is described by the direction of the twin axis (i.e. <hkl>). Inversion twinning is typically equivalent to a reflection or rotation symmetry.
Rotational twin laws are almost always 2-fold rotations, though any other permitted rotation symmetry (3-fold, 4-fold, 5-fold or 6-fold) is possible. The twin axis will be perpendicular to a lattice plane. It is possible for a rotational twin law to share the same axis as a rotational symmetry of the individual crystal if the twin law is a 2-fold rotation and the symmetry operation is a 3-fold rotation. This is the case for spinel law twinning on <111>: The spinel structure has a 3-fold rotational symmetry on <111> and spinel is commonly twinned by 2-fold rotation on <111>.
The boundary between crystal segments is called a composition surface or, if it is planar, a composition plane. The composition plane is often, though not always, parallel to the twin law plane of a reflection law. If this is the case, the twin plane is always parallel to a possible crystal face.
Common twin laws
In the isometric system, the most common types of twins are the Spinel Law (twin plane, parallel to an octahedron) <111>, where the twin axis is perpendicular to an octahedral face, and the Iron Cross <001>, which is the interpenetration of two pyritohedrons, a subtype of dodecahedron.
In the hexagonal system, calcite shows the contact twin laws {0001} and {0112}. Quartz shows the Brazil Law {1120}, and Dauphiné Law <0001>, which are penetration twins caused by transformation, and Japan Law {1122}, which is often caused by accidents during growth.
In the tetragonal system, cyclical contact twins are the most commonly observed type of twin, such as in rutile titanium dioxide and cassiterite tin oxide.
In the orthorhombic system, crystals usually twin on planes parallel to the prism face, where the most common is a {110} twin, which produces cyclical twins, such as in aragonite, chrysoberyl, and cerussite.
In the monoclinic system, twins occur most often on the planes {100} and {001} by the Manebach Law {001}, Carlsbad Law [001], Baveno Law {021} in orthoclase, and the Swallow Tail Twins (Manebach law) {001} in gypsum.
In the triclinic system, the most commonly twinned crystals are the feldspar minerals plagioclase and microcline. These minerals show the Albite and Pericline Laws.
The most common twin operations by crystal system are tabulated below. This list is not exhaustive, particularly for the crystal systems of lowest symmetry, such as the triclinic system.
Types of twinning
Simple twinned crystals may be contact twins or penetration twins. Contact twins meet on a single composition plane, often appearing as mirror images across the boundary. Plagioclase, quartz, gypsum, and spinel often exhibit contact twinning. Merohedral twinning occurs when the lattices of the contact twins superimpose in three dimensions, such as by relative rotation of one twin from the other. An example is metazeunerite. Contact twinning characteristically creates reentrant faces where faces of the crystal segments meet on the contact plane at an angle greater than 180°.
A type of twinning involving 180° relationships is called hemitropism or hemitropy.
In penetration twins the individual crystals have the appearance of passing through each other in a symmetrical manner. Orthoclase, staurolite, pyrite, and fluorite often show penetration twinning. The composition surface in penetration twins is usually irregular and extends to the center of the crystal.
Contact twinning can arise from either reflection or rotation, whereas penetration twinning is usually produced by rotation.
If several twin crystal parts are aligned by the same twin law they are referred to as multiple or repeated twins. If these multiple twins are aligned in parallel they are called polysynthetic twins. When the multiple twins are not parallel they are cyclic twins. Albite, calcite, and pyrite often show polysynthetic twinning. Closely spaced polysynthetic twinning is often observed as striations or fine parallel lines on the crystal face. Cyclic twins are caused by repeated twinning around a rotation axis. This type of twinning occurs around three, four, five, six, or eight-fold axes, the corresponding patterns are called threelings, fourlings, fivelings, sixlings, and eightlings. Sixlings are common in aragonite. Rutile, aragonite, cerussite, and chrysoberyl often exhibit cyclic twinning, typically in a radiating pattern.
For rotational twinning the relationship between the twin axis and twin plane falls into one of three types:
parallel twinning, when the twin axis and compositional plane lie parallel to each other,
normal twinning, when the twin plane and compositional plane lie normally, and
complex twinning, a combination of parallel twinning and normal twinning on one compositional plane.
Modes of formation
There are three modes of formation of twinned crystals.
Growth twins are the result of an interruption or change in the lattice during formation or growth. This may be due to a larger substituting ion, statistics as the energy difference to nucleate a new plane of atoms in a twin orientation is small, or because the twins lead to a lower energy structure.
Annealing or transformation twins are the result of a change in crystal system during cooling as one form becomes unstable and the crystal structure must re-organize or transform into another more stable form.
Deformation or gliding twins are the result of stress on the crystal after the crystal has formed. Because growth twins are formed during the initial growth of the crystal, they are described as primary, whereas transformation or deformation twins are formed in an existing crystal and are described as secondary.
Growth twinning (nanotwinning)
There are two types of twinning that can occur during growth, accidental and ones where the twinned structure has lower energy.
In accidental growth twinning an atom joins a crystal face in a less than ideal position, forming a seed for growth of a twin. The original crystal and its twin then grow together and closely resemble each other. This is characteristic enough of certain minerals to suggest that it is thermodynamically or kinetically favored under conditions of rapid growth.
Different from these are twins found in nanoparticles such as the image here, these fivefold or decahedral nanoparticles being one of the most common. These cyclic twins occur as they are lower in energy at small sizes. For the five-fold case shown, there is a disclination along the common axis which leads to an additional strain energy. Balancing this there is a reduction in the surface free energy, in large part due to more (111) surface facets. In small nanoparticles the decahedral and a more complicated Icosahedral structure (with twenty units) are lower energy, but at larger energies single crystals become lower energy. However, they do not have to transform into single crystals and can grow very large, and are known as fivelings, documented as early as 1831 by Gustav Rose; further drawings are available in the Atlas der Kristallformen, and see also the article on fivelings.
Transformation twinning
Transformation and annealing twinning takes place when a cooling crystal experiences a displacive polymorphic transition. For example, leucite has an isometric crystal structure above about , but becomes tetragonal below this temperature. Any one of the three original axes of a crystal can become the long axis when this phase change takes place. Twinning results when different parts of the crystal break their isometric symmetry along a different choice of axis. This is typically polysynthetic twinning, which enables the crystal to maintain its isometric shape by averaging out the displacement in each direction. This produces a pseudomorphic crystal that appears to have isometric symmetry. Potassium feldspar likewise experiences polysynthetic twinning as it transforms from a monoclinic structure (orthoclase) to a triclinic structure (microcline) on slow cooling.
Deformation twinning
Deformation twinning is a response to shear stress. The crystal structure is displaced along successive planes of the crystal, a process also called glide. The twinning is always reflection twinning and the glide plane is also the mirror plane. Deformation twinning can be observed in a calcite cleavage fragment by applying gentle pressure with a knife blade near an edge. This particular glide twinning, {102}, is found almost universally in deformed rock beds containing calcite.
Twinning and slip are competitive mechanisms for crystal deformation. Each mechanism is dominant in certain crystal systems and under certain conditions. In fcc metals, slip is almost always dominant because the stress required is far less than twinning stress.
Twinning can occur by cooperative displacement of atoms along the face of the twin boundary. This displacement of a large quantity of atoms simultaneously requires significant energy to perform. Therefore, the theoretical stress required to form a twin is quite high. It is believed that twinning is associated with dislocation motion on a coordinated scale, in contrast to slip, which is caused by independent glide at several locations in the crystal.
Compared to slip, twinning produces a deformation pattern that is more heterogeneous in nature. This deformation produces a local gradient across the material and near intersections between twins and grain boundaries. The deformation gradient can lead to fracture along the boundaries, particularly in bcc transition metals at low temperatures.
Of the three common crystalline structures bcc, fcc, and hcp, the hcp structure is the most likely to form deformation twins when strained, because they rarely have a sufficient number of slip systems for an arbitrary shape change. High strain rates, low stacking-fault energy and low temperatures facilitate deformation twinning.
If a metal with face-centered cubic (fcc) structure, like Al, Cu, Ag, Au, etc., is subjected to stress, it will experience twinning. The formation and migration of twin boundaries is partly responsible for ductility and malleability of fcc metals.
Twin boundaries are partly responsible for shock hardening and for many of the changes that occur in cold work of metals with limited slip systems or at very low temperatures. They also occur due to martensitic transformations: the motion of twin boundaries is responsible for the pseudoelastic and shape-memory behavior of nitinol, and their presence is partly responsible for the hardness due to quenching of steel. In certain types of high strength steels, very fine deformation twins act as primary obstacles against dislocation motion. These steels are referred to as 'TWIP' steels, where TWIP stands for twinning-induced plasticity.
Deformation twinning crystallography
Twinning is crystallographically defined by its twin plane 𝑲𝟏, the mirror plane in the twin and parent material, and 𝜼𝟏, which is the twinning shear direction. Deformation twins in Zr are generally lenticular in shape, lengthening in the 𝜼𝟏 direction and thickening along the 𝑲𝟏 plane normal.
The twin plane, shear direction, and shear plane form the basis vectors of an orthogonal set. The axis-angle misorientation relationship between the parent and twin is a rotation of angle 𝜉 about the shear plane's normal direction 𝑷.
More generally, twinning can be described as a 180° rotation about an axis (𝑲𝟏 for type I twins or 𝜼𝟏 for type II twins normal direction) , or a mirror reflection in a plane (𝑲𝟏 or 𝜼𝟏 normal plane).
In addition to a homogeneous shear, atomic shuffles are sometimes required to reform the correct crystal structure in the twinned lattice. For each twin variant, a reciprocal twin with swapped 𝑲𝟏 and 𝑲2, 𝜼𝟏 and 𝜼2 is possible, but one variant may appear more frequently in reality due to complexities with the required shuffles.
there are only two crystallographic planes in a shearing action that do not change their shape and size as a consequence of the shear. The first 𝑲𝟏 is the plane defining the upper and lower surfaces of the sheared volume. This plane contains the shear direction. The other plane, designated C. The shear direction is shown with an arrow and labelled with its customary designation 𝜼𝟏. It follows from the above that there are three ways that a crystal lattice can be sheared while still retaining its crystal structure and symmetry:
When 𝑲𝟏 is a rational plane and 𝜼2 a rational direction, a twin of the first kind
When 𝑲2 is a rational plane and 𝜼𝟏 a rational direction, a twin of the second kind, rare
When all four elements 𝑲𝟏, 𝑲2, 𝜼𝟏, and 𝜼2 are rational, a compound twin
Deformation twinning configuration
A deformation twin embryo forms in BCC metal by accumulating stacking faults, with a variant selection governed by the local stress state. Variation of the stress field close to twins inferred from HR-EBSD experimental and crystal plasticity finite element (CPFE) simulation data indicated that twins nucleate on sites with maximum strain energy density and twin resolved shear stress; thus, reducing the total elastic energy after formation. This relaxation depends on the twin thickness and is a deciding factor in the spacing between twins. Experimental and three-dimensional analysis has focussed on the (stored) strain energy density measured along a path. This highly localised stress field can provide a sufficient driving force for concurrent twin nucleation and inter/intra-granular crack nucleation.
Deformation twin growth can be perceived as a two-step process of i) thickening that is mediated by the interaction between the residual and mobile twin partials at the coherent twin-parent interface, and ii) dislocation mobility along the twin shear direction. The twin propagates when the homogeneous shear stress reaches a critical value, and a twin-parent interface advances inside the parent grain [240]. The propagating deformation twin generates a stress field due to its confinement by the surrounding parent crystal, and deformation twins develop a 3D oblate spheroid shape (which appears in 2D sections as a bi-convex lens) with a mixed coherent and non-coherent interface (Figure b).
Kannan et al. found, using in-situ ultra-high-speed optical imaging, that twin nucleation in single-crystal magnesium is stress-driven accompanied by instantaneous propagation at a speed of 1 km/s (initially) that prioritises volume lateral thickening over forward propagation, past a critical width where growth is then become faster along the shear direction. Barnett also indicated that growth is due to twin tip extension. Furthermore, elastic simulations of the local stress field surrounding the ellipsoidal twin tip find that the field can be described using its lens angle () and that the stress field magnitude increases with twin thickness.
In practice, plastic accommodation occurs in the parent crystal; thus, it also depends on the material’s yield stress, the anisotropic elastic stiffness of the parent crystal lattice, and the deformation twinning shear magnitude. This can also be accompanied by long-range diffusion of elements and elemental segregation (e.g., Cr and Co in single crystal Ni-based superalloy MD2), which occurs at the twin boundary to facilitate twin growth by lowering the critical stacking fault energy. A linear variation has been observed between twin thickness, stacking fault energy and grain size, and to a lesser degree, the stress state of the twinning grain (Schmid Factor). The twin thickness saturated once a critical residual dislocations’ density reached the coherent twin-parent crystal boundary.
Significant attention has been paid to the crystallography, morphology and macro mechanical effects of deformation twinning. Although the criterion for deformation twin growth is not entirely understood, it is a tip-controlled phenomenon linked to the interaction between the residual and mobile twin partials at the twin interface; thermodynamically, this involves the elastic energy of the strained lattice, the interface and volume free-energy of the twin, and the dissipated energy of the growth mechanism. To fully understand the interactions between microstructure (i.e., grain size, texture), temperature and strain rate on deformation twinning, it is crucial to characterise the (high) local stress and strain field associated with twin thickening and propagation. This is especially important for materials where cleavage fracture can be initiated by twinning (e.g., iron-silicon, the ferrite phase of age-hardened duplex stainless-steel, and single-crystal magnesium) as a stress-relieving mechanism.
Early studies of deformation twins arrested within grains of niobium and iron visualised the highly local strain concentration at the twin tip using an etch-pit procedure. More recently, high-resolution electron backscatter diffraction (HR-EBSD) has been used to investigate the strain 'singularity' ahead of a twin tip in hexagonal close-packed (HCP) zirconium alloy. A deformation twin in commercial purity titanium was characterised similarly and then quantified using a local Schmid factor (LSF) at the twin tip, as described in equation below.
where σ is the stress tensor, Si is the Schmid tensor, Pi is its symmetric part, di is the shear direction and ni is the shear plane normal for ith slip system. The authors concluded that conditions at the twin tip control thickening and propagation in a manner analogous to the operation of dislocation sources ahead of a crack-tip. In the analysis, a broad region of high LSF ahead of the twin tip favoured propagation, whereas a narrow region of high LSF promoted thickening. Since then, it has been argued that the LSF firmly controls the twin variant selection, as twinning has strong polarity.
The LSF novelty – compared to other criteria to describe conditions at the twin – lies in combining a geometrical criterion with the deformation field in the parent grain to provide an approximate indication of the local twin mode (i.e., thickening or propagation). However, the LSF analysis does not take advantage of the available full-field data, relies on global information on the applied stress, and does not consider the energy balance that drives twin growth. There have been few in-situ experiments to quantify the strain field ahead of a propagating deformation twin. Such observations might validate geometrical or hybrid geometrical-energy-based criteria for growth. Nanoscale testing (i.e., transmission electron microscopy) may not represent the behaviour in bulk samples due to plasticity starvation, i.e., large surface area to volume ratio, so a suitable analysis method is needed.
Lloyd described the stress concentration field ahead of the twin tip using a two-dimensional dislocation-based model within a single magnesium grain. Wang and Li, who considered microscopic phase-field (MPF) models of cracks, noted that the stress fields were similar for dislocations, deformation twinning and martensitic transformations, with differences only in the traction of the created surface, i.e., there is 100% traction recovery for dislocations and a traction-free surface for a crack. They highlighted that the stress field singularity regulates the advancement of the crack-tip and dislocations. This stress concentration can be characterised using a path-independent line integral, as shown by Eshelby for dislocations considering the contribution from the surface traction and ellipsoidal inclusions, and Rice for cracks and stress concentrations with traction-free surfaces. Furthermore, Venables noted that the oblate spheroid shape of the twin tip is the ideal example of an ellipsoid inclusion or a notch.
| Physical sciences | Crystallography | Physics |
1998254 | https://en.wikipedia.org/wiki/Granary | Granary | A granary, also known as a grain house and historically as a granarium in Latin, is a post-harvest storage building primarily for grains or seeds. Granaries are typically built above the ground to prevent spoilage and protect the stored grains or seeds from rodents, pests, floods, and adverse weather conditions. They also assist in drying the grains to prevent mold growth. Modern granaries may incorporate advanced ventilation and temperature control systems to preserve the quality of the stored grains.
Early origins
From ancient times grain has been stored in bulk. The oldest granaries yet found date back to 9500 BC and are located in the Pre-Pottery Neolithic A settlements in the Jordan Valley. The first were located in places between other buildings. However beginning around 8500 BC, they were moved inside houses, and by 7500 BC storage occurred in special rooms. The first granaries measured 3 x 3 m on the outside and had suspended floors that protected the grain from rodents and insects and provided air circulation.
These granaries are followed by those in Mehrgarh in the Indus Valley from 6000 BC. The ancient Egyptians made a practice of preserving grain in years of plenty against years of scarcity. The climate of Egypt is very dry, grain could be stored in pits for a long time without discernible loss of quality.
Historically, a silo was a pit for storing grain. It is distinct from a granary, which is an above-ground structure.
East Asia
Simple storage granaries raised on four or more posts appeared in the Yangshao culture in China and after the onset of intensive agriculture in the Korean peninsula during the Mumun pottery period (c. 1000 B.C.) as well as in the Japanese archipelago during the Final Jōmon/Early Yayoi periods (c. 800 B.C.). In the archaeological vernacular of Northeast Asia, these features are lumped with those that may have also functioned as residences and together are called 'raised floor buildings'.
China built an elaborate system designed to minimize famine deaths. The system was destroyed in the Taiping Rebellion of the 1850s.
Southeast Asia
In vernacular architecture of Indonesian archipelago granaries are made of wood and bamboo materials and most of them are built and raised on four or more posts to avoid rodents and pests. Examples of Indonesian granaries styles are the Sundanese leuit and Minang rangkiang.
Great Britain
In the South Hams in southwest Great Britain, small granaries were built on mushroom-shaped stumps called staddle stones. They were built of timber-frame construction and often had slate roofs. Larger ones were similar to linhays but with the upper floor enclosed. Access to the first floor was usually via a stone staircase on the outside wall.
Towards the close of the 19th century, warehouses specially intended for holding grain began to multiply in Great Britain. There are climatic difficulties in the way of storing grain in Great Britain on a large scale, but these difficulties have been largely overcome.
Moisture control
Grain must be kept away from moisture for as long as possible to preserve it in good condition and prevent mold growth. Newly harvested grain brought into a granary tends to contain excess moisture, which encourages mold growth leading to fermentation and heating, both of which are undesirable and affect quality. Fermentation generally spoils grain and may cause chemical changes that create poisonous mycotoxins.
One traditional remedy is to spread the grain in thin layers on a floor, where it is turned to aerate it thoroughly. Once the grain is sufficiently dry it can be transferred to a granary for storage. Today, this can be done using a mechanical grain auger to move grain from one granary to another.
In modern silos, grain is typically force-aerated in situ or circulated through external grain drying equipment.
Modern
Modern grain farming operations often use manufactured steel granaries to store grain on-site until it can be trucked to major storage facilities in anticipation of shipping. The large mechanized facilities, particularly seen in Russia and North America are known as grain elevators.
Examples
| Technology | Buildings and infrastructure | null |
1998561 | https://en.wikipedia.org/wiki/Scorpius%20X-1 | Scorpius X-1 | Scorpius X-1 is an X-ray source located roughly 9000 light years away in the constellation Scorpius. Scorpius X-1 was the first extrasolar X-ray source discovered, and, aside from the Sun, it is the strongest apparent non-transient source of X-rays in the sky. The X-ray flux varies day-to-day, and is associated with an optically visible star, V818 Scorpii, that has an apparent magnitude which fluctuates between 12-13.
Discovery and early study
The possible existence of cosmic soft X-rays was first proposed by Bruno Rossi, MIT professor and board chairman of American Science and Engineering in Cambridge, Massachusetts to Martin Annis, president of AS&E. Following his urging, the company obtained a contract from the United States Air Force to explore the lunar surface prior to the launch of astronauts to the Moon, and incidentally to perhaps see galactic sources of X-rays.
Subsequently, Scorpius X-1 was discovered in 1962 by a team, under Riccardo Giacconi, who launched an Aerobee 150 sounding rocket carrying a highly sensitive soft X-ray detector designed by Frank Paolini. The rocket trajectory was slightly off course but still detected a significant emission of soft X-rays that were not coming from the Moon. Thus fortuitously, and as first pointed out by Frank Paolini, Scorpius X-1 became the first X-ray source discovered outside the Solar System. The angular resolution of the detector did not initially allow the position of Scorpius X-1 to be accurately determined. This led to suggestions that the source might be located near the Galactic Center, but it was eventually realized that it lies in the constellation Scorpius. As the first discovered X-ray source in Scorpius, it received the designation Scorpius X-1.
The Aerobee 150 rocket launched on June 12, 1962, detected the first X-rays from another celestial source (Scorpius X-1) at J1950 RA Dec . Sco X-1 is a LMXB in which the visual counterpart is V818 Scorpii.
Although the above reference indicates the rocket launch was on June 12, 1962, other sources indicate the actual launch was at 06:59:00 UTC on June 19, 1962.
Historical footnote:
"The instrumentation had been designed for an attempt to observe X-rays from the moon and was not equipped with collimation to restrict the field of view narrowly. As a result, the signal was very broad, and accurate definition of the size and position of the source was not possible. A similar experiment was repeated in October 1962 when the Galactic Center was below the horizon and the strong source was not present. A third attempt, in June 1963, verified the results of the June 1962 flight." The Galactic Center is < 20° RA and < 20° Dec from Sco X-1, the two X-ray sources are separated by ~20° of arc and may not have been resolvable in the June 1962 flight.
Scorpius XR-1 has been observed at J1950 RA Dec .
In 1967 (before the discovery of pulsars), Iosif Shklovsky examined X-ray and optical observations of Scorpius X-1 and correctly concluded that the radiation comes from a neutron star accreting matter from a companion.
Characteristics
Its X-ray output is 2.3×1031 W, about 60,000 times the total luminosity of the Sun. Scorpius X-1 shows regular variations of up to 1 magnitude in its intensity, with a period of around 18.9 hours. The source varies irregularly in optical wavelengths as well, but these changes are not correlated with the X-ray variations. Scorpius X-1 itself is a neutron star whose intense gravity draws material off its companion into an accretion disk, where it ultimately falls onto the surface, releasing a tremendous amount of energy. As this stellar material accelerates in Scorpius X-1's gravitational field, X-rays are emitted. The measured luminosity for Scorpius X-1 is consistent with a neutron star which is accreting matter at its Eddington limit.
This system is classified as a low-mass X-ray binary; the neutron star is roughly 1.4 solar masses, while the donor star is only 0.42 solar masses. The two stars were probably not born together; recent research suggests that the binary may have been formed by a close encounter inside a globular cluster.
| Physical sciences | Notable stars | Astronomy |
1999394 | https://en.wikipedia.org/wiki/Sand%20shark | Sand shark | Sand sharks, also known as sand tiger sharks, gray nurse sharks or ragged tooth sharks, are mackerel sharks of the family Odontaspididae. They are found worldwide in temperate and tropical waters. The three species are in two genera.
Description
The body tends to be brown with dark markings in the upper half. These markings disappear as they mature. Their needle-like teeth are highly adapted for impaling fish, their main prey. Their teeth are long, narrow, and very sharp with smooth edges, with one and on occasion two smaller cusplets on either side. Sand sharks have a large second dorsal fin. The sand shark can grow up to long, and most adults can weigh around . The average lifespan of both sexes is only about 7 years, though they may live longer in captivity.
Location and origins
The name sand shark comes from their tendency to migrate toward shoreline habitats, and they are often seen swimming around the ocean floor in the surf zone; at times, they come very close to shore. They are often found in warm or temperate waters throughout the world's oceans, except the eastern Pacific. They also frequent the Mediterranean and Adriatic Seas at depths from and sometimes more.
Behavior
The sand shark has a unique hunting strategy. It is able to gulp air from above the surface and collect the air in its stomach. This enables it to become buoyant and approach its prey virtually motionless. During the day, the sand shark stays mostly inactive, but at night, it becomes active and resumes hunting activities. Its staple is small fish, but it eats crustaceans and squid, as well. It occasionally hunts in shivers (groups), and has even been known to attack full fishing nets.
Reproduction
Sand sharks only develop two embryos, one in each uterus. The largest and strongest embryos consume their siblings in the womb (intrauterine cannibalism) before each surviving pup is born. It has one of the lowest reproduction rates of all sharks and is susceptible to even minimal population pressure, so it is listed as vulnerable and is protected in much of its range.
Attacks on people
Sand sharks are not known to attack humans. If a person were to provoke a sand shark, it may retaliate defensively. Sand sharks are generally not aggressive, but harass divers who are spearfishing. In North America, wreck divers regularly visit the World War II shipwrecks to dive with the sharks that make the wrecks their home.
Conservation
A recent report from the PEW Charitable Trusts suggests a new management approach used for large mammals that have suffered population declines could hold promise for sharks. Because of the life-history characteristics of sharks, conventional fisheries management approaches, such as reaching maximum sustainable yield, may not be sufficient to rebuild depleted shark populations. Some of the more stringent approaches used to reverse declines in large mammals may be appropriate for sharks, including prohibitions on the retention of the most vulnerable species and regulation of international trade.
Species
The family contains three extant species, in two genera, as well as many extinct species in several genera. Recent mitochondrial DNA analysis of extant members has found the two extant members do not actually form a monophyletic clade. This family is therefore polyphyletic and in need of revision.
Genus Carcharias Rafinesque, 1810
Carcharias taurus Rafinesque, 1810 (sand tiger shark)
Genus Odontaspis Agassiz 1838
Odontaspis ferox A. Risso, 1810 (smalltooth sand tiger)
Odontaspis noronhai Maul, 1955 (bigeye sand tiger)
Subfamily Odontaspinae † Herman, 1975
Genus Striatolamia † Glikman, 1964
Genus Carcharoides † Ameghino, 1901
Genus Parodontaspis † White, 1931
Genus Priodontaspis † Ameghino, 1901
Genus Pseudoisurus † Glikman, 1957
Genus Synodontaspis † White, 1931
Subfamily Johnlonginae † Shimada, 2015
Genus Johnlongia †
Genus Pseudomegachasma † (Shimada, 2015)
| Biology and health sciences | Sharks | Animals |
1999621 | https://en.wikipedia.org/wiki/Styrax | Styrax | Styrax (common names storax or snowbell) is a genus of about 130 species of large shrubs or small trees in the family Styracaceae, mostly native to warm temperate to tropical regions of the Northern Hemisphere, with the majority in eastern and southeastern Asia, but also crossing the equator in South America. The resin obtained from the tree is called benzoin or storax (not to be confused with the Liquidambar storax balsam), often used as a vanilla-like component in perfumery.
The genus Pamphilia, sometimes regarded as distinct, is now included within Styrax based on analysis of morphological and DNA sequence data. The spicebush (Lindera benzoin) is a different plant, in the family Lauraceae.
Styrax trees grow to 2–14 m tall, and have alternate, deciduous or evergreen simple ovate leaves 1–18 cm long and 2–10 cm broad. The flowers are pendulous, with a white 5–10-lobed corolla, produced 3–30 together on open or dense panicles 5–25 cm long. The fruit is an oblong dry drupe, smooth and lacking ribs or narrow wings, unlike the fruit of the related snowdrop trees (Halesia) and epaulette trees (Pterostyrax).
Uses
Uses of resin
Benzoin resin, a dried exudation from pierced bark, is currently produced from various Styrax species native to Sumatra, Java, and Thailand. Commonly traded are the resins of S. tonkinensis (Siam benzoin), S. benzoin (Sumatra benzoin), and S. benzoides. The name benzoin is probably derived from Arabic lubān jāwī (لبان جاوي, "Javan frankincense); compare the obsolete terms gum benjamin and benjoin. This incidentally shows that the Arabs were aware of the origin of these resins, and that by the late Middle Ages at latest international trade in them was probably of major importance.
The chemical benzoin (2-hydroxy-2-phenylacetophenone), despite the apparent similarity of the name, is not contained in benzoin resin in measurable quantities. However, benzoin resin does contain small amounts of the hydrocarbon styrene, named however for Levant storax (from Liquidambar orientalis), from which it was first isolated, and not for the genus Styrax itself; industrially produced styrene is now used to produce polystyrene plastics, including Styrofoam.
History of sources
Since Antiquity, storax resin has been used in perfumes, certain types of incense, and medicines.
There is some degree of uncertainty as to exactly what resin old sources refer to. Turkish sweetgum (Liquidambar orientalis) is a quite unrelated tree in the family Altingiaceae that produces a similar resin traded in modern times as storax or as Levant storax, like the resins of other sweetgums, and a number of confusing variations thereupon. Turkish sweetgum is a relict species that occurs only in a small area in SW Turkey (and not in the Levant at all); presumably, quite some of the "storax resin" of the Ancient Greek and the Ancient Roman sources was from this sweetgum, rather than a Styrax, although at least during the former era genuine Styrax resin, probably from S. officinalis, was imported in quantity from the Near East by Phoenician merchants, and Herodotus of Halicarnassus in the 5th century BC indicates that different kinds of storax were traded.
The nataf (נטף) of the incense sacred to Yahweh, mentioned in the Book of Exodus, is loosely translated by the Greek term staktē (στακτή, AMP: ), or an unspecific "gum resin" or similar term (NIV: ). Nataf may have meant the resin of Styrax officinalis or of some other plant, perhaps Turkish sweetgum, which is unlikely to have been imported in quantity into the Near East.
Since the Middle Ages, Southeast Asian benzoin resins became increasingly available; today there is little international trade in S. officinalis resin and little production of Turkish sweetgum resin due to that species' decline in numbers.
Use as incense
Storax incense is used in the Middle East and adjacent regions as an air freshener. This was adopted in the European Papier d'Arménie. Storax resin from southern Arabian species was burned during frankincense (Boswellia resin) harvesting; it was said to drive away snakes:
"[The Arabians] gather frankincense by burning that storax which Phoenicians carry to Hellas; they burn this and so get the frankincense; for the spice-bearing trees are guarded by small winged snakes of varied color, many around each tree; these are the snakes that attack Egypt. Nothing except the smoke of storax will drive them away from the trees."
Medical uses
There has been little dedicated research into the medical properties of storax resin, but it has been used for long, and apparently with favorable results. It was important in Islamic medicine; Avicenna (Ibn Sina, ابن سینا) discusses S. officinalis it in his Al-Qanun fi al-Tibb (القانون في الطب, The Law of Medicine). He indicates that storax resin mixed with other antibiotic substances and hardening material gives a good dental restorative material. Benzoin resin is a component of the "Theriaca Andromachi Senioris", a Venice treacle recipe in the 1686 d'Amsterdammer Apotheek.
Tincture of benzoin is benzoin resin dissolved in alcohol. This and its numerous derived versions like lait virginal and friar's balsam were highly esteemed in 19th-century European cosmetics and other household purposes; they apparently had antibacterial properties. Today tincture of benzoin is most often used in first aid for small injuries, as it acts as a disinfectant and local anesthetic and seems to promote healing. Benzoin resin and its derivatives are also used as additives in cigarettes.
The antibiotic activity of benzoin resin seems mostly due to its abundant benzoic acid and benzoic acid esters, which were named after the resin; other less well known secondary compounds such as lignans like pinoresinol are likely significant too.
Horticultural uses
Several species of storax are popular ornamental trees in parks and gardens, especially S. japonicus and its cultivars such as 'Emerald Pagoda', and Styrax obassia.
Uses of wood
The wood of larger species is suitable for fine handicrafts. That of egonoki (エゴノキ, S. japonicus) is used to build kokyū (胡弓), the Japanese bowed instrument.
Ecology and conservation
The resin of Styrax acts to kill wound pathogens and deter herbivores. Consequently, for example, few Lepidoptera caterpillars eat storax compared to other plants. Those of the two-barred flasher (Astraptes fulgerator) were recorded on S. argenteus, but they do not seem to use it on a regular basis.
Some storax species have declined in numbers due to unsustainable logging and habitat degradation. While most of these are classified as vulnerable (VU) by the IUCN, only four trees of the nearly extinct palo de jazmin (S. portoricensis) are known to survive at a single location. Although legally protected, this species could be wiped out by a single hurricane.
Selected species
Styrax agrestis – China
Styrax americanus – SE USA
Styrax argenteus – N & S America
Styrax argentifolius – China
Styrax bashanensis – China
Styrax benzoides – Thailand, S China
Styrax benzoin – Sumatra
Styrax calvescens – China
Styrax camporum – Brazil, Bolivia, Paraguay
Styrax chinensis – China
Styrax chrysocalyx – Brazil
Styrax chrysocarpus – China
Styrax confusus – China
Styrax cordatus – Peru and Ecuador
Styrax crotonoides – Malaysia
Styrax dasyanthus – central China
Styrax faberi – China
Styrax ferrugineus – Brazil, Bolivia, Paraguay
Styrax formosanus – China
Styrax foveolaria – Peru and Ecuador
Styrax fraserensis – Malaysia
Styrax grandiflorus – China
Styrax grandifolius – SE USA
Styrax hainanensis – S China
Styrax hemsleyanus – China
Styrax hookeri – Himalaya
Styrax huanus – China
Styrax jaliscana – Mexico
Styrax japonicus – Japan
Styrax limpritchii – SW China (Yunnan)
Styrax litseoides – Vietnam
Styrax macranthus – China
Styrax macrocarpus – China
Styrax martii – Brazil
Styrax obassia – Japan, China
Styrax odoratissimus – China
Styrax officinalis – SE Europe, SW Asia
Styrax pentlandianus – Bolivia, Peru, Ecuador, Colombia
Styrax perkinsiae – China
Styrax peruvianus – Costa Rica, Panama, Colombia, Ecuador, Peru
Styrax philadelphoides – China
Styrax platanifolius – Texas, NE Mexico
Styrax pohlii – Suriname, Brazil, Peru, Bolivia
Styrax portoricensis – Puerto Rico
Styrax redivivus – California
Styrax roseus – China
Styrax rugosus – China
Styrax schweliense – W China
Styrax serrulatus – Himalaya, SW China
Styrax shiraianum – Japan
Styrax suberifolius – China
Styrax supaii – China
Styrax tomentosus – Colombia, Ecuador and Peru
Styrax tonkinensis – SE Asia
Styrax veitchiorum – China
Styrax vilcabambae – Peru
Styrax wilsonii – W China
Styrax wuyuanensis – China
Styrax zhejiangensis – China
| Biology and health sciences | Ericales | Plants |
32826316 | https://en.wikipedia.org/wiki/Pinterest | Pinterest | Pinterest is an American social media service for publishing and discovery of information in the form of pinboards. This includes recipes, home, style, motivation, and inspiration on the Internet using image sharing. Pinterest, Inc. was founded by Ben Silbermann, Paul Sciarra, and Evan Sharp, and is headquartered in San Francisco.
History
Pinterest emerged from an earlier app created by Ben Silbermann and Paul Sciarra called Tote which served as a virtual replacement for paper catalogs. Tote struggled as a business, significantly due to difficulties with mobile payments. At the time, mobile payment technology was not sophisticated enough to enable easy on-the-go transactions, inhibiting users from making many purchases via the app. Tote users however were amassing large collections of favorite items and sharing them with other users. The behavior struck a chord with Silbermann, and he shifted the company to building Pinterest, which allowed users to create collections of a variety of items and share them with each other.
The development of Pinterest began in December 2009, and the site launched the prototype as a closed beta in March 2010. Nine months after the launch, the website had 10,000 users. Silbermann said he wrote to the first 5,000 users, offering his phone number and even meeting with some of them. The launch of an iPhone app in early March 2011 brought in more downloads than expected. This was followed by an iPad app and Pinterest Mobile, a version of the website for non-iPhone users. Silbermann and a few programmers operated the site out of a small apartment until mid-2011.
Pinterest grew rapidly during this period. On August 10, 2011, Time magazine listed Pinterest in its "50 Best Websites of 2011" article. In December 2011, the site became one of the top 10 largest social network services, according to Hitwise data, with 11 million total visits per week. Pinterest won the Best New Startup of 2011 at the TechCrunch Crunchies Awards. For January 2012, comScore reported the site had 11.7 million unique U.S. visitors, making it the fastest site ever to break through the 10 million unique visitor mark. At the 2012 Webby Awards, Pinterest won Best Social Media App and People's Voice Award for best functioning visual design.
On March 23, 2012, Pinterest unveiled updated terms of service that eliminated the policy that gave it the right to sell its users' content. On August 10, 2012, Pinterest altered its policy so that a request or an invitation was no longer required to join the site. In October 2012, Pinterest launched business accounts allowing businesses to either convert their existing personal accounts into business accounts or start from scratch.
In April 2017, Pinterest removed its post "liking" feature as it seemed redundant to "boards", which are user collections of posts. Users' existing indexes of liked posts were converted into a collection ("board") named as such.
Although starting out as a "social network" with boards, in later years the company has put increasing emphasis in visual search and e-commerce, such as shopping catalogs.
In February 2019, The Wall Street Journal stated that Pinterest secretly filed for an initial public offering (IPO) of stock. The total valuation of the company at the time reached $12 billion. They went public on April 18, 2019, at $19 per share, closing the day at $24.40 per share.
For 2020, Pinterest reported advertising revenue of $1.7 billion, an increase of 48% from 2019. On March 3, 2021, Pinterest announced Pinterest Premiere, a video ads product "which will appear in people's feeds, targeted to their interests and other characteristics". Later in April, chief financial officer Todd Morgenfeld announced plans to spend more money on marketing in order to offset a potential slowdown in activity as the United States economy reopened with more people getting vaccinated for COVID-19.
On October 20, 2021, Bloomberg reported that PayPal is interested in acquiring Pinterest, with a potential price of around $70 a share. PayPal's board and management decided later that same week to back away from a potential deal.
In December 2021, Pinterest acquired the editing and video creation app Vochi. Following this, In May 2022, it was announced that Pinterest released a new video streaming app “Pinterest TV studio”. The app is aimed at allowing users to live-stream on its platform and use different devices for different angles while live-streaming on the Pinterest platform.
On June 28, 2022, Pinterest announced that co-founder, CEO and President, Ben Silbermann would transition to the newly created role of Executive Chairman, and online commerce expert Bill Ready will become Chief Executive Officer and a member of the Board of Directors.
On August 1, 2022, Pinterest quietly launched a new app named Shuffles, which allows users to build collages out of photos from Pinterest's library or those that are uploaded by the user.
In January 2023 at CES, Pinterest announced its partnership with LiveRamp, a data enablement platform to create data 'clean rooms' for selected advertisers on the platform. These 'clean rooms' allow Pinterest's ad partners to utilize first-party data for personalized ads without having to share the data with Pinterest. With data privacy a large concern for online platforms and their users, this partnership is an effort to stimulate ad business on the platform while keeping its user's data safe and in compliance with new data collection regulations. The first advertiser to pilot this feature will be grocery retailer, Albertsons with a winter healthy eating campaign.
Features and content
The creators behind Pinterest summarized the service as a "catalogue of ideas" that inspires users to "go out and do that thing", although that it is not an image-based "social network". It also has a very large fashion profile. In later years, Pinterest has also been described as a "visual search engine".
Pinterest consists mainly of "pins" and "boards", where a pin is an image that has been linked from a website or uploaded. Pins saved from one user's board can be saved to someone else's board, a process known as "repinning". Boards are collections of pins dedicated to a theme. Boards with multiple ideas can have different sections that further contain multiple pins. Users can follow and unfollow other users as well as boards, which would fill the "home feed".
Content can also be found outside Pinterest and similarly uploaded to a board via the "Save" button, which can be downloaded to the bookmark bar on a web browser, or be implemented by a webmaster directly on the website. It was originally called the "Pin it" button, but it was renamed in 2016 to "Save" to try to make the site more intuitive to new users.
In August 2016, Pinterest launched a video player that lets users and brands upload and store clips of any length straight to the site.
Exploring
The home feed is a collection of Pins from the users, boards, and topics followed, as well as a few promoted pins and pins Pinterest has picked. On the main Pinterest page, a "pin feed" appears, displaying the chronological activity from the Pinterest boards that a user follows.
In October 2013, Pinterest began displaying advertisements in the form of "Promoted Pins". Promoted Pins are based on an individual user's interests, things done on Pinterest, or a result of visiting an advertiser's site or app.
In 2015, Pinterest implemented a feature that allows users to search with images instead of words.
In March 2020, Pinterest introduced the "Today" tab on the home feed which shows trending pins.
In October 2022, Pinterest announced that its video-focused “Idea Pins” feature has included the ability to add popular tracks from top artists from licensing deals with Warner Music Group, Warner Chappell Music, Merlin, and BMG.
Visual search
In 2017, Pinterest introduced a "visual search" function that allows users to search for elements in images (existing pins, existing parts of a photo, or new photos) and guide users to suggested similar content within Pinterest's database. The tools powered by artificial intelligence are called Pinterest Lens, Shop the Look, and Instant Ideas.
Shopping and catalogs
The platform has drawn businesses, especially retailers, to create pages aimed at promoting their companies online as a "virtual storefront".
In 2013, Pinterest introduced a new tool called "Rich Pins", to enhance the customer experience when browsing through pins made by companies. Business pages can include various data, topics, and information such as prices of products, ratings of movies or ingredients for recipes.
In June 2015, Pinterest unveiled "buyable pins" that allows users to purchase things directly from Pinterest. In October 2018, the buyable pins feature was replaced by "Product Pins"
In March 2019, Pinterest added product catalogs and personalized shopping recommendations with the "more from [brand]" option, showcasing a range of product Pins from the same business.
Pinterest Analytics
Pinterest Analytics is a service that generates comprehensive statistics on a specific website's traffic, commonly used by marketers. Impressions, Engagements, Pin clicks, Outbound clicks, and Saves are some aspects of user data that Pinterest Analytics provides. It also collects data that depicts the percentage of change within a specific time to determine if a product is more popular on a specific day during the week or is slowly becoming unpopular. This data can help marketing agencies. The "Most Clicked" tab in Pinterest Analytics demonstrates products that are more likely to sell. Through the access of Pinterest Analytics, companies receive data via API.
Creator Fund
Pinterest has the conception of creating a digital and visual bookmarking platform for displaying personal identities and tastes. The Pinterest interface emphasizes the content of pins and links while the users are the media through which the content is communicated, spreading among users through pinning and repinning. The individualized formulation of the app makes it not very social; instead, it serves its pursuit to discover, develop, and refine personal interests. Pinterest faces a decline in the period of platforms like TikTok.
In 2021, the site announced its new Creator Fund, with the stated aim of trying to support creators and their ability to monetize their efforts to preserve engagement and interactions on the platform. The program's initial launch increased creators' overall monthly views by 72%. Pinterest will invest $1.2 million in underrepresented creators via cash grants, ad credits, and equipment.
Teen Safety Features
Pinterest profiles that claimed to be aged 16 and under were defaulted to private, and all of their followers were removed. Moreover, new affordances have been created surrounding group board invites and messaging within the application, restricting these functions to mutual followers. This rollout is following the review by NBC News by Stephen Sauer of the nonprofit Canadian centre for Child Protection, who found that "his homepage 'almost immediately' filled with images of children often dressed in similarly revealing attire, several of which had received sexually suggestive comments". In order to further the agenda of protecting young users, Pinterest has added more affordances for reporting 'spam', 'inappropriate cover images', and options to 'call-out when content may involve a minor'.
Usage
Pinterest is a free website that requires registration to use. The service is currently accessible through a web browser, and apps for iOS, Android, and Windows 10 PCs.
In February 2013, Reuters and ComScore stated that Pinterest had 48.7 million users globally, and a study released in July 2013 by French social media agency Semiocast revealed the website had 70 million users worldwide. In October 2016, the company had 150 million monthly active users (70 million in the U.S. and 80 million outside it), rising to 175 million monthly active users by April 2017 and 250 million in September 2018. As of July 2020, there were over 400 million monthly active users. In April 2023, Pinterest reported 463 million monthly active users, which suggests that 7.4% of the world's population over age 13 use Pinterest.
Around 2020, Pinterest was thought to flood search results in Google Images. In 2022 Google claimed to have performed changes to increase "diversity" in the search results. Pinterest claimed that, as a consequence of Google's changes of November 2021, "U.S. monthly active users coming to Pinterest from the web, desktop and mobile web declined around 30% year over year".
Demographics
Pinterest has largely appealed to women, especially with its early user base. A 2020 report found that over 60% of the global users are women. Although men have not been a primary audience on Pinterest, their usage reportedly increased 48%. In terms of age distribution, users between the ages of 18 and 25 have grown twice as fast as those over the age of 25. However, both users between the ages of 18 and 25 and users between the ages of 25 and 40 have been driving the growth of Pinterest.
In science
Data from Pinterest has been used for research in different areas. For example, it is possible to find patterns of activity that attract the attention of audience and content reposting, including the extent to which users specialize in particular topics, and homophily among users. Another work focused on studying the characteristics, manifestations and overall effects of user behaviors from various aspects, as well as correlations between neighboring users and the topology of the network structure. There is also a study that was based on Pinterest and proposed a novel pinboard recommendation system for Twitter users.
Corporate affairs
Pinterest, Inc. is headquartered in San Francisco, California. Originally it was based in Palo Alto before moving in 2012.
In early 2011, the company secured a US$10 million Series A financing led by Jeremy Levine and Sarah Tavel of Bessemer Venture Partners. In October 2011, the company secured US$27 million in funding from Andreessen Horowitz, which valued the company at US$200 million. Co-founder Paul Sciarra left his position at Pinterest in April 2012 for a consulting job as entrepreneur in residence at Andreessen Horowitz.
On May 17, 2012, Japanese electronic commerce company Rakuten announced it was leading a $100 million investment in Pinterest, alongside investors including Andreessen Horowitz, Bessemer Venture Partners, and FirstMark Capital, based on a valuation of $1.5 billion.
On September 20, 2012, Pinterest announced the hiring of its new head of engineering, Jon Jenkins. Jenkins came from Amazon, where he spent eight years as an engineering lead and was also a director of developer tools, platform analysis and website platform.
In late October 2013, Pinterest secured a $225 million round of equity funding that valued the website at $3.8 billion.
In 2014, Pinterest generated its first revenue, when it began charging advertisers. An analyst at Wedbush Securities estimated that ads could generate up to $500 million in 2016. In 2015, investors valued Pinterest, Inc. at $11 billion, making it a "unicorn" (a start-up with a valuation exceeding $1 billion). , the company was valued at $12 billion.
In June 2017, Pinterest raised $150 million from a group of existing investors.
In August 2020, Pinterest paid $89.5 million to cancel a large office space lease on a to-be-completed complex in San Francisco's SoMa area, near its current headquarters.
In April 2023, Pinterest announced a partnership with Amazon to show third-party advertisements on the website, which will allow users to be redirect to Amazon to make purchases.
Acquisitions
In March 2013, Pinterest acquired Livestar. The terms were not disclosed. In early October 2013, Pinterest acquired Hackermeter. The company's co-founders, Lucas Baker and Frost Li, joined Pinterest as engineers.
In April 2015, Pinterest acquired the team from Hike Labs, which had been developing a mobile publishing application called Drafty.
In May 2016, Pinterest acquired mobile deep linking startup URX to help accelerate its content understanding efforts. The URX team's expertise in mobile content discovery and recommendation would prove critical to helping Pinterest understand its corpus of over 100 billion pins, to better recommend them to its users.
On August 23, 2016, Pinterest announced that it would be acquiring the team behind Instapaper, which will continue operating as a separate app. The Instapaper team will both work on the core Pinterest experience and updating Instapaper.
On March 8, 2017, Pinterest said it had acquired Jelly Industries, a small search-engine company founded by Biz Stone.
In December 2021, Pinterest announced the acquisition of the Vochi app.
In June 2022, Pinterest announced its definitive agreement to acquire San Francisco based AI-driven fashion shopping platform, The Yes. Pinterest announced that it had closed the acquisition on June 10, 2022.
Criticism
Copyrighted content
Pinterest has a notification system that copyright holders can use to request that content be removed from the site. The Digital Millennium Copyright Act (DMCA) safe harbor status of Pinterest has been questioned given that it actively promotes its users to copy to Pinterest, for their perpetual use, any image on the Internet. Pinterest users cannot claim safe harbor status and as such are exposed to possible legal action for pinning copyright material. Pinterest allows users to transfer information; intellectual property rights come to play.
A "nopin" HTML meta tag was released by Pinterest on February 20, 2012, to allow websites to opt out of their images being pinned. On February 24, 2012, Flickr implemented the code to allow users to opt out.
Pinterest released a statement in March 2012 saying it believed it was protected by the DMCA's safe harbor provisions.
In early May 2012, the site added automatic attribution of authors on images originating from Flickr, Behance, YouTube and Vimeo. Automatic attribution was also added for Pins from sites mirroring content on Flickr. At the same time, Flickr added a Pin shortcut to its share option menu to users who have not opted out of sharing their images.
Content creators on sites such as iStock have expressed concern over their work being reused on Pinterest without permission. Getty Images said that it was aware of Pinterest's copyright issues and was in discussion with them.
Legal status
In February 2012, photographer and lawyer Kirsten Kowalski wrote a blog post explaining how her interpretation of copyright law led her to delete all her infringing pins. The post contributed to scrutiny over Pinterest's legal status. The post went viral and reached founder Ben Silbermann who contacted Kowalski to discuss making the website more compliant with the law.
Terms of service
Pinterest's earlier terms of service asserted the right to "sell" user content. In a March 2012 article in Scientific American, illustrator Kalliopi Monoyios alleged that "Pinterest's terms of service have been garnering a lot of criticism for stating in no uncertain terms that anything you 'pin' to their site belongs to them. Completely. Wholly. Forever and for always." At the time, Pinterest's terms of service stated that:
According to Monoyios, Pinterest's claim to a broad license to sell user content potentially undermined artists' ability to monetize their own work. Another Scientific American blogger claimed that this provision contradicted another line in the terms of service, that "Cold Brew Labs does not claim any ownership rights in any such Member Content".
Several days later, Pinterest unveiled updated terms of service that, once implemented in April, ended the site's previous claims of ownership of posted images. "Selling content was never our intention", said the company in a blog post.
Censorship
In its 2019 "Who Has Your Back?" report, the Electronic Frontier Foundation gave Pinterest a three (out of six) star rating, highlighting improvements in the company's transparency reports about government takedown notices, but criticizing the lack of a clear commitment to notify users about content removals and account suspensions.
In March 2017, Chinese authorities blocked Pinterest without explanation. The block was imposed during the annual National People's Congress, a politically sensitive period in the country. While Pinterest is not known for its political content, experts identified the ban as consistent with Chinese government efforts to use website blocks and the "Great Firewall" as an industrial policy tool to promote Chinese tech companies (e.g., Baidu, Youku, Weibo, and Renren) by censoring foreign tech companies. Huaban, Duitang and many other websites bear similarities to Pinterest.
Internet service providers in India had blocked Pinterest following a Madras High Court order in July 2016 to block a list of around 225 "rogue websites indulging in online piracy and infringement of copyright". The block was temporary.
Content policies and user bans
In October 2012, Pinterest added a new feature allowing users to report others for negative and offensive activity or block other users if they do not want to view their content, a bid that the company said aimed to keep the site "positive and respectful."
In December 2018, Pinterest began to take steps to block health misinformation from its recommendations engine, and blocked various searches, content, and user accounts that related to, or promoted, unproved and disproven cancer treatments. The company said it also blocked multiple accounts that linked to external websites that sold supplements and other products that were not scientifically validated. In January 2019, Pinterest stopped returning search results relating to vaccines, in an effort to somehow slow the increase of anti-vaccination content on the platform. Prior to the measure, the company said that the majority of vaccination-related images shared on the platform were anti-vaccination, contradicting the scientific research establishing the safety of vaccines.
In June 2019, anti-abortion group Live Action was banned from Pinterest; the company said the permanent suspension was imposed for spreading "harmful misinformation, [which] includes medical misinformation and conspiracies that turn individuals and facilities into targets for harassment or violence."
In December 2019, following a campaign from the activist group Color of Change, Pinterest announced that it would restrict content that advertises wedding events on former slave plantations.
Culture of discrimination
In 2020, two former Pinterest employees, Ifeoma Ozoma and Aerica Shimizu Banks went public about their experience at Pinterest. Both women recounted experiences of discrimination at work, including racist comments, unequal pay, and punishment for speaking out. Additionally, Ozoma claims that the company failed to protect her when personal information was shared with hate sites by a colleague of hers. In response, Pinterest released an apology statement and CEO Ben Silbermann sent an email to all employees pushing the company to do better.
In August 2020, dozens of Pinterest staff participated in a virtual walkout in support of two former colleagues who publicly accused the company of racism and gender discrimination.
In December 2020, Pinterest agreed to pay its former Chief Operating Officer $20 million+ to settle a lawsuit alleging discrimination.
In November 2021, Pinterest settled a lawsuit that alleged racial and gender discrimination. The company agreed to spend $50 million on improving its diversity and to release former employees from non-disclosure agreements. The settlement was in regard to Ozoma and Banks's accusations of June 2020.
| Technology | Social network and blogging | null |
30231169 | https://en.wikipedia.org/wiki/Interbreeding%20between%20archaic%20and%20modern%20humans | Interbreeding between archaic and modern humans | Interbreeding between archaic and modern humans occurred during the Middle Paleolithic and early Upper Paleolithic. The interbreeding happened in several independent events that included Neanderthals and Denisovans, as well as several unidentified hominins.
In Europe, Asia and North Africa, interbreeding between archaic humans and modern humans took place several times. The introgression events into modern humans are estimated to have happened about 47,000–65,000 years ago with Neanderthals and about 44,000–54,000 years ago with Denisovans.
Neanderthal-derived DNA has been found in the genomes of most or possibly all contemporary populations, varying noticeably by region. It accounts for 1–4% of modern genomes for people outside Sub-Saharan Africa, although estimates vary, and either none or up to 0.3% for those in Sub-Saharan Africa. Cushitic and Semitic speaking populations from the Horn of Africa (such as Ethiopians), who derive a large portion of their ancestry from West Eurasians, have ~1% Neanderthal-derived DNA.
Neanderthal-derived DNA is highest in East Asians, intermediate in Europeans, and lower in Southeast Asians. According to some research, it is also lower in Melanesians and Polynesians compared to both East Asians and Europeans. However, other research finds higher Neanderthal admixture in Melanesians, as well as in Native Americans, than in Europeans (though not higher than in East Asians).
Denisovan-derived ancestry is largely absent from modern populations in Africa, Western Asia and Europe. The highest rates, by far, of Denisovan admixture have been found in Oceanian and some Southeast Asian populations. An estimated 4–6% of the genome of modern Melanesians is derived from Denisovans, but the highest amounts detected thus far are found in the Negrito populations of the Philippines. While some Southeast Asian Negrito populations carry Denisovan admixture, others, such as the Andamanese, have none. In addition, low traces of Denisovan-derived ancestry have been found in mainland Asia, with an elevated Denisovan ancestry in South Asian populations compared to other mainland populations.
In Africa, archaic alleles consistent with several independent admixture events in the subcontinent have been found. It is currently unknown who these archaic African hominins were. A 2020 paper found that "despite their very low levels or absence of archaic ancestry, African populations share many Neanderthal and Denisovan variants that are absent from Eurasia, reflecting how a larger proportion of the ancestral human variation has been maintained in Africa."
A 2016 paper in the journal Evolutionary Biology argued that introgression of DNA from other lineages enabled humanity to migrate to, and succeed in, numerous new environments, with the resulting hybridization being an essential force in the emergence of modern humans.
In December 2023, scientists reported that genes inherited by modern humans from Neanderthals and Denisovans may biologically influence the daily routine of modern humans.
Neanderthals
Genetics
Proportion of admixture
On 7 May 2010, following the genome sequencing of three Vindija Neanderthals, a draft sequence of the Neanderthal genome was published and revealed that Neanderthals shared more alleles with Eurasian populations (e.g. French, Han Chinese, and Papua New Guinean) than with sub-Saharan African populations (e.g. Yoruba and San). According to the authors Green et al. (2010), the observed excess of genetic similarity is best explained by recent gene flow from Neanderthals to modern humans after the migration out of Africa. They estimated the proportion of Neanderthal-derived ancestry to be 1–4% of the Eurasian genome. Durand et al. (2011) estimated 1–6% Neanderthal ancestry in non-Africans. Prüfer et al. (2013) estimated the proportion to be 1.5–2.1% for non-Africans. Lohse and Frantz (2014) infer a higher rate of 3.4–7.3% in Eurasia. In 2017, Prüfer et al. revised their estimate to 1.8–2.6% for non-Africans outside Oceania.
According to a later study by Chen et al. (2020), Africans (specifically, the 1000 Genomes African populations) also have Neanderthal admixture, with this Neanderthal admixture in African individuals accounting for 17 megabases, which is 0.3% of their genome. According to the authors, Africans gained their Neanderthal admixture predominantly from a back-migration by peoples (modern humans carrying Neanderthal admixture) that had diverged from ancestral Europeans (postdating the split between East Asians and Europeans). This back-migration is proposed to have happened about 20,000 years ago. However, some scientists, such as geneticist David Reich, have doubts about how extensive the flow of DNA back to Africa would have been, finding the signal of Neanderthal admixture "really weak".
Introgressed genome
It has been found that 50% of the Neanderthal genome is present among people in India, and 41% has been found in Icelanders.
Previously it was found that about 20% of the Neanderthal genome was found in modern Eurasians, but the figure was also estimated at a third. A 2023 study found an introgression from modern humans to Neanderthals around 250,000 years ago, and estimated that roughly 6% of the Altai Neanderthal genome was inherited from modern humans.
Subpopulation admixture rate
A higher Neanderthal admixture was found in East Asians than in Europeans, which is estimated to be about 20% more introgression into East Asians. This could possibly be explained by the occurrence of further admixture events in the early ancestors of East Asians after the separation of Europeans and East Asians, dilution of Neanderthal ancestry in Europeans by populations with low Neanderthal ancestry from later migrations, or reduced efficacy of purifying selection in the ancestors of East Asians, due smaller effective population sizes as they migrated to East Asia. Studies simulating admixture models indicate that a reduced efficacy of purifying selection against Neanderthal alleles in East Asians could not account for the greater proportion of Neanderthal ancestry of East Asians, thus favoring more-complex models involving additional pulses of admixture between Neanderthals and the ancestors of East Asians. Such models show a pulse to ancestral Eurasians, followed by separation and an additional pulse to ancestral East Asians. It is observed that there is a small but significant variation of Neanderthal admixture rates within European populations, but no significant variation within East Asian populations. Prüfer et al. (2017) remarked that East Asians carry more Neanderthal DNA (2.3–2.6%) than Western Eurasians (1.8–2.4%).
It was later determined by Chen et al. (2020) that East Asians have 8% more Neanderthal ancestry, revised from the previous reports of 20% more Neanderthal ancestry, compared to Europeans. This stems from the fact that Neanderthal ancestry shared with Africans had been masked, because Africans were thought to have no Neanderthal admixture and were therefore used as reference samples. Thus, any overlap in Neanderthal admixture with Africans resulted in an underestimation of Neanderthal admixture in non-Africans and especially in Europeans. The authors give a single pulse of Neanderthal admixture after the out-of-Africa dispersal as the most parsimonious explanation for the enrichment in East Asians, but they add that variation in Neanderthal ancestry may also be attributed to dilution to account for the now-more-modest differences found. As a proportion of the total amount of Neanderthal sequence for each population, 7.2% of the sequence in Europeans is shared exclusively with Africans, while 2% of the sequence in East Asians is shared exclusively with Africans.
Genomic analysis suggests that there is a global division in Neanderthal introgression between sub-Saharan African populations and other modern human groups (including North Africans) rather than between African and non-African populations. North African groups share a similar excess of derived alleles with Neanderthals as do non-African populations, whereas sub-Saharan African groups are the only modern human populations that generally did not experience Neanderthal admixture. The Neanderthal genetic signal among North African populations was found to vary depending on the relative quantity of North African, European, Near Eastern and sub-Saharan ancestry. Using F4 ancestry ratio statistical analysis, the Neanderthal inferred admixture was observed to be: highest among the North African populations with highest North African ancestry such as Tunisian Berbers, where it was at the same level or even higher than that of Eurasian populations (100–138%); high among North African populations carrying greater European or Near Eastern admixture, such as groups in North Morocco and Egypt (~60–70%); and lowest among North African populations with greater Sub-Saharan admixture, such as in South Morocco (20%). Quinto et al. (2012) therefore postulate that the presence of this Neanderthal genetic signal in Africa is not due to recent gene flow from Near Eastern or European populations since it is higher among populations bearing indigenous pre-Neolithic North African ancestry. Low but significant rates of Neanderthal admixture has also been observed for the Maasai of East Africa. After identifying African and non-African ancestry among the Maasai, it can be concluded that recent non-African modern human (post-Neanderthal) gene flow was the source of the contribution since around an estimated 30% of the Maasai genome can be traced to non-African introgression from about 100 generations ago.
Distance to lineages
Presenting a high-quality genome sequence of a female Altai Neanderthal, it has been found that the Neanderthal component in non-African modern humans is more related to the Mezmaiskaya Neanderthal (North Caucasus) than to the Altai Neanderthal (Siberia) or the Vindija Neanderthals (Croatia). By high-coverage sequencing the genome of a 50,000-year-old female Vindija Neanderthal fragment, it was later found that the Vindija and Mezmaiskaya Neanderthals did not seem to differ in the extent of their allele-sharing with modern humans. In this case, it was also found that the Neanderthal component in non-African modern humans is more closely related to the Vindija and Mezmaiskaya Neanderthals than to the Altai Neanderthal. These results suggest that a majority of the admixture into modern humans came from Neanderthal populations that had diverged (about 80–100kya) from the Vindija and Mezmaiskaya Neanderthal lineages before the latter two diverged from each other.
Analysis of chromosome 21 of the Altai, El Sidrón (Spain), and Vindija Neanderthals indicates that of these three lineages, only the El Sidrón and Vindija Neanderthals display significant rates of gene flow (0.3–2.6%) into modern humans, suggesting that the El Sidrón and Vindija Neanderthals are more closely related than the Altai Neanderthal to the Neanderthals that interbred with modern humans about 47,000–65,000 years ago. Conversely, significant rates of modern human gene flow into Neanderthals occurred—of the three examined lineages—for only the Altai Neanderthal (0.1–2.1%), suggesting that modern human gene flow into Neanderthals mainly took place after the separation of the Altai Neanderthals from the El Sidrón and Vindija Neanderthals that occurred roughly 110,000 years ago. The findings show that the source of modern human gene flow into Neanderthals originated from a population of early modern humans from about 100,000 years ago, predating the out-of-Africa migration of the modern human ancestors of present-day non-Africans.
Mitochondrial DNA and Y chromosome
No evidence of Neanderthal mitochondrial DNA has been found in modern humans. This suggests that successful Neanderthal admixture happened in pairings with Neanderthal males and modern human females. Possible hypotheses are that Neanderthal mitochondrial DNA had detrimental mutations that led to the extinction of carriers, that the hybrid offspring of Neanderthal mothers were raised in Neanderthal groups and became extinct with them, or that female Neanderthals and male Sapiens did not produce fertile offspring. However, the hypothesized incompatibility between Neanderthals and modern humans is contested by findings that suggest that the Y chromosome of Neanderthals was replaced by an extinct lineage of the modern human Y chromosome, which introgressed into Neanderthals between 100,000 and 370,000 years ago. Furthermore, the study concludes that the replacement of the Y chromosomes and mitochondrial DNA in Neanderthals after gene flow from modern humans is highly plausible given the increased genetic load in Neanderthals relative to modern humans.
As shown in an interbreeding model produced by Neves and Serva (2012), the Neanderthal admixture in modern humans may have been caused by a very low rate of interbreeding between modern humans and Neanderthals, with the exchange of one pair of individuals between the two populations in about every 77 generations. This low rate of interbreeding would account for the absence of Neanderthal mitochondrial DNA from the modern human gene pool as found in earlier studies, as the model estimates a probability of only 7% for a Neanderthal origin of both mitochondrial DNA and Y chromosome in modern humans.
Reduced contribution
There are large genomic regions with strongly reduced Neanderthal contribution in modern humans due to negative selection, partly caused by hybrid male infertility. These regions were most-pronounced on the X chromosome, with fivefold lower Neanderthal ancestry compared to autosomes. They also contained relatively high numbers of genes specific to testes. This means that modern humans have relatively few Neanderthal genes that are located on the X chromosome or expressed in the testes, suggesting male infertility as a probable cause. It may be partly affected by hemizygosity of X chromosome genes in males.
Deserts of Neanderthal sequences may also be caused by genetic drift involving intense bottlenecks in the modern human population and background selection as a result of strong selection against deleterious Neanderthal alleles. The overlap of many deserts of Neanderthal and Denisovan sequences suggests that repeated loss of archaic DNA occur at specific loci.
It has also been shown that Neanderthal ancestry has been selected against in conserved biological pathways, such as RNA processing.
Consistent with the hypothesis that purifying selection has reduced Neanderthal contribution in present-day modern human genomes, Upper Paleolithic Eurasian modern humans (such as the Tianyuan modern human) carry more Neanderthal DNA (about 4–5%) than present-day Eurasian modern humans (about 1–2%).
Rates of selection against Neanderthal sequences varied for European and Asian populations.
Changes in modern humans
In Eurasia, modern humans have adaptive sequences introgressed from archaic humans, which provided a source of advantageous genetic variants that are adapted to local environments and a reservoir for additional genetic variation. Adaptive introgression from Neanderthals has targeted genes involved with keratin filaments, sugar metabolism, muscle contraction, body fat distribution, enamel thickness, and oocyte meiosis, as well as brain size and functioning. There are signals of positive selection, as the result of adaptation to diverse habitats, in genes involved with variation in skin pigmentation and hair morphology. In the immune system, introgressed variants have heavily contributed to the diversity of immune genes, of which there's an enrichment of introgressed alleles that suggest a strong positive selection.
Genes affecting keratin were found to have been introgressed from Neanderthals into modern humans (shown in East Asians and Europeans), suggesting that these genes gave a morphological adaptation in skin and hair to modern humans to cope with non-African environments. This is likewise for several genes involved in medical-relevant phenotypes, such as those affecting systemic lupus erythematosus, primary biliary cirrhosis, Crohn's disease, optic disk size, smoking behavior, interleukin 18 levels, and diabetes mellitus type 2.
Researchers found Neanderthal introgression of 18 genes—several of which are related to UV-light adaptation—within the chromosome 3p21.31 region (HYAL region) of East Asians. The introgressive haplotypes were positively selected in only East Asian populations, rising steadily from 45,000 years BP until a sudden increase of growth rate around 5,000 to 3,500 years BP. They occur at very high frequencies among East Asian populations in contrast to other Eurasian populations (e.g. European and South Asian populations). The findings also suggests that this Neanderthal introgression occurred within the ancestral population shared by East Asians and Native Americans.
Evans et al. (2006) had previously suggested that a group of alleles collectively known as haplogroup D of microcephalin, a critical regulatory gene for brain volume, originated from an archaic human population. The results show that haplogroup D introgressed 37,000 years ago (based on the coalescence age of derived D alleles) into modern humans from an archaic human population that separated 1.1 million years ago (based on the separation time between D and non-D alleles), consistent with the period when Neanderthals and modern humans co-existed and diverged respectively. The high frequency of the D haplogroup (70%) suggest that it was positively selected for in modern humans. The distribution of the D allele of microcephalin is high outside Africa but low in sub-Saharan Africa, which further suggest that the admixture event happened in archaic Eurasian populations. This distribution difference between Africa and Eurasia suggests that the D allele originated from Neanderthals according to Lari et al. (2010), but they found that a Neanderthal individual from the Mezzena Rockshelter (Monti Lessini, Italy) was homozygous for an ancestral allele of microcephalin, thus providing no support that Neanderthals contributed the D allele to modern humans and also not excluding the possibility of a Neanderthal origin of the D allele. Green et al. (2010), having analyzed the Vindija Neanderthals, also could not confirm a Neanderthal origin of haplogroup D of the microcephalin gene.
It has been found that HLA-A*02, A*26/*66, B*07, B*51, C*07:02, and C*16:02 of the immune system were contributed from Neanderthals to modern humans. After migrating out of Africa, modern humans encountered and interbred with archaic humans, which was advantageous for modern humans in rapidly restoring HLA diversity and acquiring new HLA variants that are better adapted to local pathogens.
It is found that introgressed Neanderthal genes exhibit cis-regulatory effects in modern humans, contributing to the genomic complexity and phenotype variation of modern humans. Looking at heterozygous individuals (carrying both Neanderthal and modern human versions of a gene), the allele-specific expression of introgressed Neanderthal alleles was found to be significantly lower in the brain and testes relative to other tissues. In the brain, this was most pronounced at the cerebellum and basal ganglia. This downregulation suggests that modern humans and Neanderthals possibly experienced a relative higher rate of divergence in these specific tissues.
Furthermore, correlating the genotypes of introgressed Neanderthal alleles with the expression of nearby genes, it is found that archaic alleles contribute proportionally more to variation in expression than nonarchaic alleles. Neanderthal alleles affect expression of the immune genes OAS1/2/3 and TLR1/6/10, which can be specific to cell type and is influenced by environmental stimuli.
Studying the high-coverage female Vindija Neanderthal genome, Prüfer et al. (2017) identified several Neanderthal-derived gene variants, including those that affect levels of LDL cholesterol and vitamin D, and that influence eating disorders, visceral fat accumulation, rheumatoid arthritis, schizophrenia, as well as responses to antipsychotic drugs.
Examining European modern humans in regards to the Altai Neanderthal genome in high-coverage, results show that Neanderthal admixture is associated with several changes in cranium and underlying brain morphology, suggesting changes in neurological function through Neanderthal-derived genetic variation. Neanderthal admixture is associated with an expansion of the posterolateral area of the modern human skull, extending from the occipital and inferior parietal bones to bilateral temporal locales. In regards to modern human brain morphology, Neanderthal admixture is positively correlated with an increase in sulcal depth for the right intraparietal sulcus and an increase in cortical complexity for the early visual cortex of the left hemisphere. Neanderthal admixture is also positively correlated with an increase in white and gray matter volume localized to the right parietal region adjacent to the right intraparietal sulcus. In the area overlapping the primary visual cortex gyrification in the left hemisphere, Neanderthal admixture is positively correlated with gray matter volume. The results also show evidence for a negative correlation between Neanderthal admixture and white matter volume in the orbitofrontal cortex.
In Papuans, Neanderthal genetic variants are found in highest frequency in genes expressed in the brain, whereas Denisovan DNA has the highest frequency in genes expressed in bones and other tissues.
A Neanderthal allele inherited by modern humans, SNP rs3917862, is with associated with hypercoagulability. This can be harmful, but women lacking the allele are 0.1% more likely to die in childbirth.
In December 2023, scientists reported that genes inherited by modern humans from Neanderthals and Denisovans may biologically influence the daily routine of modern humans.
Population substructure theory
Although less parsimonious than recent gene flow, the observation may have been due to ancient population sub-structure in Africa, causing incomplete genetic homogenization within modern humans when Neanderthals diverged while early ancestors of Eurasians were still more closely related to Neanderthals than those of Africans were to Neanderthals. On the basis of allele frequency spectrum, it was shown that the recent admixture model had the best fit to the results while the ancient population sub-structure model had no fit—demonstrating that the best model was a recent admixture event that was preceded by a bottleneck event among modern humans – thus confirming recent admixture as the most parsimonious and plausible explanation for the observed excess of genetic similarities between modern non-African humans and Neanderthals. On the basis of linkage disequilibrium patterns, a recent admixture event is likewise confirmed by the data. From the extent of linkage disequilibrium, it was estimated that the last Neanderthal gene flow into early ancestors of Europeans occurred 47,000–65,000 years BP. In conjunction with archaeological and fossil evidence, the gene flow is thought likely to have occurred somewhere in Western Eurasia, possibly the Middle East. Through another approach—using one genome each of a Neanderthal, Eurasian, African, and chimpanzee (outgroup), and dividing it into non-recombining short sequence blocks—to estimate genome-wide maximum-likelihood under different models, an ancient population sub-structure in Africa was ruled out and a Neanderthal admixture event was confirmed.
Morphology
The early Upper Paleolithic burial remains of a modern human child from Abrigo do Lagar Velho (Portugal) features traits that indicate Neanderthal interbreeding with modern humans dispersing into Iberia. Considering the dating of the burial remains (24,500 years BP) and the persistence of Neanderthal traits long after the transitional period from a Neanderthal to a modern human population in Iberia (28,000–30,000 years BP), the child may have been a descendant of an already heavily admixed population.
The remains of an early Upper Paleolithic modern human from Peștera Muierilor (Romania) of 35,000 years BP shows a morphological pattern of European early modern humans, but possesses archaic or Neanderthal features, suggesting European early modern humans interbreeding with Neanderthals. These features include a large interorbital breadth, a relatively flat superciliary arches, a prominent occipital bun, an asymmetrical and shallow mandibular notch shape, a high mandibular coronoid processus, the relative perpendicular mandibular condyle to notch crest position, and a narrow scapular glenoid fossa.
The early modern human Oase 1 mandible from Peștera cu Oase (Romania) of 34,000–36,000 14C years BP presents a mosaic of modern, archaic, and possible Neanderthal features. It displays a lingual bridging of the mandibular foramen, not present in earlier humans except Neanderthals of the late Middle and Late Pleistocene, thus suggesting affinity with Neanderthals. Concluding from the Oase 1 mandible, there was apparently a significant craniofacial change of early modern humans from at least Europe, possibly due to some degree of admixture with Neanderthals.
The earliest (before about 33 ka BP) European modern humans and the subsequent (Middle Upper Paleolithic) Gravettians, falling anatomically largely in line with the earliest (Middle Paleolithic) African modern humans, also show traits that are distinctively Neanderthal, suggesting that a solely Middle Paleolithic modern human ancestry was unlikely for European early modern humans.
Manot 1, a partial calvarium of a modern human that was recently discovered at the Manot Cave (Western Galilee, Israel) and dated to 54.7±5.5 kyr BP, represents the first fossil evidence from the period when modern humans successfully migrated out of Africa and colonized Eurasia. It also provides the first fossil evidence that modern humans inhabited the southern Levant during the Middle to Upper Palaeolithic interface, contemporaneously with the Neanderthals and close to the probable interbreeding event. The morphological features suggest that the Manot population may be closely related to or may have given rise to the first modern humans who later successfully colonized Europe to establish early Upper Palaeolithic populations.
History
The interbreeding has been discussed ever since the discovery of Neanderthal remains in the 19th century, though earlier writers believed that Neanderthals were a direct ancestor of modern humans. Thomas Huxley suggested that many Europeans bore traces of Neanderthal ancestry, but associated Neanderthal characteristics with primitivism, writing that since they "belong to a stage in the development of the human species, antecedent to the differentiation of any of the existing races, we may expect to find them in the lowest of these races, all over the world, and in the early stages of all races".
Until the early 1950s, most scholars thought Neanderthals were not in the ancestry of living humans. Nevertheless, Hans Peder Steensby proposed interbreeding in 1907 in the article Race studies in Denmark. He strongly emphasised that all living humans are of mixed origins. He held that this would best fit observations, and challenged the widespread idea that Neanderthals were ape-like or inferior. Basing his argument primarily on cranial data, he noted that the Danes, like the Frisians and the Dutch, exhibit some Neanderthaloid characteristics, and felt it was reasonable to "assume something was inherited" and that Neanderthals "are among our ancestors".
Carleton Stevens Coon in 1962 found it likely, based upon evidence from cranial data and material culture, that Neanderthal and Upper Paleolithic peoples either interbred or that the newcomers reworked Neanderthal implements "into their own kind of tools".
By the early 2000s, the majority of scholars supported the Out of Africa hypothesis, according to which anatomically modern humans left Africa about 50,000 years ago and replaced Neanderthals with little or no interbreeding.
Yet some scholars still argued for hybridisation with Neanderthals. The most vocal proponent of the hybridisation hypothesis was Erik Trinkaus of Washington University in St. Louis. Trinkaus claimed various fossils as products of hybridised populations, including the skeleton of a child found at Lagar Velho in Portugal and the Peștera Muierii skeletons from Romania.
Denisovans
Genetics
Proportion of admixture
It has been shown that Melanesians (e.g. Papua New Guinean and Bougainville Islander) share relatively more alleles with Denisovans when compared to other Eurasian-derived populations and Africans. It is estimated that 4% to 6% of the genome in Melanesians derives from Denisovans, while no Eurasians or Africans displayed contributions of the Denisovan genes. It has been observed that Denisovans contributed genes to Melanesians but not to East Asians, indicating that there was interaction between the early ancestors of Melanesians with Denisovans but that this interaction did not take place in the regions near southern Siberia, where as-of-yet the only Denisovan remains have been found. In addition, Aboriginal Australians also show a relative increased allele sharing with Denisovans, compared to Eurasians and African populations, consistent with the hypothesis of increased admixture between Denisovans and Melanesians.
Reich et al. (2011) produced evidence that the highest presence of Denisovan admixture is in Oceanian populations, followed by many Southeast Asian populations, and none in East Asian populations. There is significant Denisovan genetic material in eastern Southeast Asian and Oceanian populations (e.g. Aboriginal Australians, Near Oceanians, Polynesians, Fijians, eastern Indonesians, Philippine Mamanwa and Manobo), but not in certain western and continental Southeast Asian populations (e.g. western Indonesians, Malaysian Jehai, Andaman Onge, and mainland Asians), indicating that the Denisovan admixture event happened in Southeast Asia itself rather than mainland Eurasia. The observation of high Denisovan admixture in Oceania and the lack thereof in mainland Asia suggests that early modern humans and Denisovans had interbred east of the Wallace Line that divides Southeast Asia according to Cooper and Stringer (2013).
Skoglund and Jakobsson (2011) observed that particularly Oceanians, followed by Southeast Asians populations, have a high Denisovans admixture relative to other populations. Furthermore, they found possible low traces of Denisovan admixture in East Asians and no Denisovan admixture in Native Americans. In contrast, Prüfer et al. (2013) found that mainland Asian and Native American populations may have a 0.2% Denisovan contribution, which is about twenty-five times lower than Oceanian populations. The manner of gene flow to these populations remains unknown. However, Wall et al. (2013) stated that they found no evidence for Denisovan admixture in East Asians.
Findings indicate that the Denisovan gene flow event happened to the common ancestors of Aboriginal Filipinos, Aboriginal Australians, and New Guineans. New Guineans and Australians have similar rates of Denisovan admixture, indicating that interbreeding took place prior to their common ancestors' entry into Sahul (Pleistocene New Guinea and Australia), at least 44,000 years ago. It has also been observed that the fraction of Near Oceanian ancestry in Southeast Asians is proportional to the Denisovan admixture, except in the Philippines where there is a higher proportional Denisovan admixture to Near Oceanian ancestry. Reich et al. (2011) suggested a possible model of an early eastward migration wave of modern humans, some who were Philippine/New Guinean/Australian common ancestors that interbred with Denisovans, respectively followed by divergence of the Philippine early ancestors, interbreeding between the New Guinean and Australian early ancestors with a part of the same early-migration population that did not experience Denisovan gene flow, and interbreeding between the Philippine early ancestors with a part of the population from a much-later eastward migration wave (the other part of the migrating population would become East Asians).
Finding components of Denisovan introgression with differing relatedness to the sequenced Denisovan, Browning et al. (2018) suggested that at least two separate episodes of Denisovan admixture has occurred. Specifically, introgression from two distinct Denisovan populations is observed in East Asians (e.g. Japanese and Han Chinese), whereas South Asians (e.g. Telugu and Punjabi) and Oceanians (e.g. Papuans) display introgression from one Denisovan population.
Exploring derived alleles from Denisovans, Sankararaman et al. (2016) estimated that the date of Denisovan admixture was 44,000–54,000 years ago. They also determined that the Denisovan admixture was the greatest in Oceanian populations compared to other populations with observed Denisovan ancestry (i.e. America, Central Asia, East Asia, and South Asia). The researchers also made the surprising finding that South Asian populations display an elevated Denisovan admixture (when compared to other non-Oceanian populations with Denisovan ancestry), albeit the highest estimate (which are found in Sherpas) is still ten times lower than in Papuans. They suggest two possible explanations: There was a single Denisovan introgression event that was followed by dilution to different extents or at least three distinct pulses of Denisovan introgressions must have occurred.
A study in 2021 analyzing archaic ancestry in 118 Philippine ethnic groups discovered an independent admixture event into Philippine Negritos from Denisovans. The Ayta Magbukon in particular were found to possess the highest level of Denisovan ancestry in the world, with ~30%–40% more than even that found in Australians and Papuans (Australo-Melanesians), suggesting that distinct Islander Denisovan populations existed in the Philippines which admixed with modern humans after their arrival.
It has been shown that Eurasians have some but significantly lesser archaic-derived genetic material that overlaps with Denisovans, stemming from the fact that Denisovans are related to Neanderthals—who contributed to the Eurasian gene pool—rather than from interbreeding of Denisovans with the early ancestors of those Eurasians.
The skeletal remains of an early modern human from the Tianyuan cave (near Zhoukoudian, China) of 40,000 years BP showed a Neanderthal contribution within the range of today's Eurasian modern humans, but it had no discernible Denisovan contribution. It is a distant relative to the ancestors of many Asian and Native American populations, but post-dated the divergence between Asians and Europeans. The lack of a Denisovan component in the Tianyuan individual suggests that the genetic contribution had been always scarce in the mainland.
Reduced contribution
There are large genomic regions devoid of Denisovan-derived ancestry, partly explained by infertility of male hybrids, as suggested by the lower proportion of Denisovan-derived ancestry on X chromosomes and in genes that are expressed in the testes of modern humans.
Changes in modern humans
Exploring the immune system's HLA alleles, it has been suggested that HLA-B*73 introgressed from Denisovans into modern humans in western Asia due to the distribution pattern and divergence of HLA-B*73 from other HLA alleles. Even though HLA-B*73 is not present in the sequenced Denisovan genome, HLA-B*73 was shown to be closely associated to the Denisovan-derived HLA-C*15:05 from the linkage disequilibrium. From phylogenetic analysis, however, it has been concluded that it is highly likely that HLA-B*73 was ancestral.
The Denisovan's two HLA-A (A*02 and A*11) and two HLA-C (C*15 and C*12:02) allotypes correspond to common alleles in modern humans, whereas one of the Denisovan's HLA-B allotype corresponds to a rare recombinant allele and the other is absent in modern humans. It is thought that these must have been contributed from Denisovans to modern humans, because it is unlikely to have been preserved independently in both for so long due to HLA alleles' high mutation rate.
Tibetan people received an advantageous EGLN1 and EPAS1 gene variant, associated with hemoglobin concentration and response to hypoxia, for life at high altitudes from the Denisovans. The ancestral variant of EPAS1 upregulates hemoglobin levels to compensate for low oxygen levels—such as at high altitudes—but this also has the maladaption of increasing blood viscosity. The Denisovan-derived variant on the other hand limits this increase of hemoglobin levels, thus resulting in a better altitude adaption. The Denisovan-derived EPAS1 gene variant is common in Tibetans and was positively selected in their ancestors after they colonized the Tibetan plateau.
Archaic African hominins
Rapid decay of fossils in Sub-Saharan African environments makes it currently unfeasible to compare modern human admixture with reference samples of archaic Sub-Saharan African hominins.
Ancient DNA Data from a ~4,500 BP Ethiopian highland individual, and from Southern (~2,300–1,300 BP), and Eastern and South-Central Africa (~8,100–400 BP) has clarified that some West Africa populations have small amounts of excess alleles best explained by an archaic source in West Africans that is not included in the pre-agricultural Eastern African hunter-gatherers, Southern African hunter-gatherer populations, or the genetic gradation between them. The West African groups carrying the archaic DNA include Yoruba from coastal Nigeria and Mende from Sierra Leon indicating that the ancient DNA was acquired long before the spread of agriculture and likely well before the Holocene (starting 11,600 BP), Such an archaic lineage must have separated before the divergence of San ancestors, which is estimated to have begun on the order of 200–300 thousand years ago.
The hypothesis that there has been archaic line in the ancestry of present-day Africans that originated before the San, Pygmies and East African hunter gatherers (and the Eurasians) is supported by a line of evidence independent from the Skoglund findings based on long haplotypes with deep divergences from other human haplotypes including Lachance et al.(2012), Hammer et al., 2011, and Plagnol and Wall (2006).
In the archaic DNA differences found by Hammer, et al., the pygmies (of Central Africa) are grouped with the San (of Southern Africa) in contrast to the Yoruba (of West Africa). Further clarification of the presence of archaic DNA in current West African populations with the extraction and sequencing of DNA from 4 fossils found at Shum Laka in Cameroon dating from 8,000 to 3,000 BP. These individuals were found to derive most of their DNA from Central African hunter gatherers (Pygmy ancestors) and did not share the archaic DNA found in the Yoruba and Mande. The pattern of differences between Eastern, Central and Southern hunter gatherers when compared to the West African groups which had been found by Hammer was confirmed. In a second study Lipson et al. (2020) studied DNA extracted from 6 additional Eastern and Southcentral African fossils from the last 18,000 years. It was determined that their genetic origins could be accounted for by DNA contributions from Southern, Central and Eastern hunter gatherers, and that none of them had the archaic DNA found in the Yoruba.
According to a study published in 2020, there are indications that 2% to 19% (or about ≃6.6 and ≃7.0%) of the DNA of four West African populations may have come from an unknown archaic hominin which split from the ancestor of humans and Neanderthals between 360 kya to 1.02 mya. However, in contrast to the studies of Skoglund and Lipson with ancient African DNA, the study also finds that at least part of this proposed archaic admixture is also present in Eurasians/non-Africans, and that the admixture event or events range from 0 to 124 ka B.P, which includes the period before the Out-of-Africa migration and prior to the African/Eurasian split (thus affecting in part the common ancestors of both Africans and Eurasians/non-Africans). Another recent study, which discovered substantial amounts of previously undescribed human genetic variation, also found ancestral genetic variation in Africans that predates modern humans and was lost in most non-Africans.
Archaic hominins in Eurasia
Hominins' presence in Eurasia began at least 2 million years BP. Genetic evidence shows that thousands of years later when lineages of Neandertals and Denisovans started to expand into Eurasia, the continent was still inhabited by descendants of these archaic hominins, and their genetic admixture made its way into genome of Neanderthals and Denisovans and later indirectly into modern humans.
Genetic studies show two major events of genetic admixture from superarchaics, suggesting that in the late middle Pleistocene, Eurasia was inhabited by at least two separate populations of ancient hominins.
Roger et al. (2020) describes an event of admixture that occurred soon after Neandersovans (common ancestor of Neanderthals and Denisovans) started to expand into Eurasia. They met a lineage of superarchaic hominins that had been separated from African homo lineages since at least 2 Ma ago.
Previous studies identified more recent event of admixture. About 350,000 years ago a genome of an "erectus-like" creature was injected into the Denisovan lineage. With the separation time of about 2 Ma ago and interbreeding that happened 350 ka ago, the two populations involved were more distantly related than any pair of human populations previously known to interbreed.
Related studies
In 2019, scientists discovered evidence, based on genetics studies using artificial intelligence (AI), that suggests the existence of an unknown human ancestor species, not Neanderthal or Denisovan, in the genome of modern humans.
| Biology and health sciences | Homo | Biology |
3738572 | https://en.wikipedia.org/wiki/Combustion%20analysis | Combustion analysis | Combustion analysis is a method used in both organic chemistry and analytical chemistry to determine the elemental composition (more precisely empirical formula) of a pure organic compound by combusting the sample under conditions where the resulting combustion products can be quantitatively analyzed. Once the number of moles of each combustion product has been determined the empirical formula or a partial empirical formula of
the original compound can be calculated.
Applications for combustion analysis involve only the elements of carbon (C), hydrogen (H), nitrogen (N), and sulfur (S) as combustion of materials containing them convert these elements to their oxidized form (CO2, H2O, NO or NO2, and SO2) under high temperature high oxygen conditions. Notable interests for these elements involve measuring total nitrogen in food or feed to determine protein percentage, measuring sulfur in petroleum products, or measuring total organic carbon (TOC) in water.
History
The method was invented by Joseph Louis Gay-Lussac. Justus von Liebig studied the method while working with
Gay-Lussac between 1822 and 1824 and improved the method in the following years to a level that it could be used as standard procedure for organic analysis.
Combustion train
A combustion train is an analytical tool for the determination of elemental composition of a chemical compound. With knowledge of elemental composition a chemical formula can be derived. The combustion train allows the determination of carbon and hydrogen in a succession of steps:
combustion of the sample at high temperatures with Copper(II) oxide as the oxidizing agent,
collection of the resulting gas in a hygroscopic agent (magnesium perchlorate or calcium chloride) to trap generated water,
collection of the remainder gas in a strong base (for instance potassium hydroxide) to trap generated carbon dioxide.
Analytical determination of the amounts of water and carbon dioxide produced from a known amount of sample gives the empirical formula. For every hydrogen atom in the compound 1/2 equivalent of water is produced, and for every carbon atom in the compound 1 equivalent of carbon dioxide is produced.
Nowadays, modern instruments are sufficiently automated to be able to do these analyses routinely. Samples required are also extremely small — 0.5 mg of sample can be sufficient to give satisfactory CHN analysis.
CHN analyzer
A CHN analyzer (also known as a carbon hydrogen and nitrogen analyzer) is a scientific instrument which is used to measure carbon, hydrogen and nitrogen elemental concentrations in a given sample with accuracy and precision. Sample sizes are most often just a few milligrams, but may differ depending on system. For some sample matrices larger mass is preferred due to sample heterogeneity. These analysers are capable of handling a wide variety of sample types, including solids, liquids, volatile and viscous samples, in the fields of pharmaceuticals, polymers, chemicals, environment, food and energy.
This instrument calculates the percentages of elemental concentrations based on the Dumas method, using flash combustion of the sample to cause an instantaneous oxidization into simple compounds which are then detected with thermal conductivity detection or infrared spectroscopy. Separation of interference is done by chemical reagents.
Modern methods
The water vapor, carbon dioxide and other products can be separated via gas chromatography and analysed via a thermal conductivity detector.
| Physical sciences | Basics_2 | Chemistry |
3739627 | https://en.wikipedia.org/wiki/Yellow-lipped%20sea%20krait | Yellow-lipped sea krait | The yellow-lipped sea krait (Laticauda colubrina), also known as the banded sea krait or colubrine sea krait, is a species of venomous snake found in tropical Indo-Pacific oceanic waters. The snake has distinctive black stripes and a yellow snout, with a paddle-like tail for use in swimming.
It spends much of its time under water to hunt, but returns to land to digest, rest, and reproduce. It has very potent neurotoxic venom, which it uses to prey on eels and small fish. Because of its affinity to land, the yellow-lipped sea krait often encounters humans, but the snake is not aggressive and only attacks when feeling threatened.
Description
The head of a yellow-lipped sea krait is black, with lateral nostrils and an undivided rostral scale. The upper lip and snout are characteristically colored yellow, and the yellow color extends backward on each side of the head above the eye to the temporal scales.
The body of the snake is subcylindrical, and is taller than it is wide. Its upper surface is typically a shade of blueish gray, while the belly is yellowish, with wide ventral scales that stretch from a third to more than half of the width of the body. Black rings of about uniform width are present throughout the length of the snake, but the rings narrow or are interrupted at the belly. The midbody is covered with 21 to 25 longitudinal rows of imbricated (overlapping) dorsal scales. The dorsal and lateral scales can be used to differentiate between this species and the similar yellow-lipped New Caledonian sea krait, which typically has fewer rows of scales and scales that narrow or fail to meet (versus the yellow-lipped sea krait's ventrally meeting dark bands). The tail of the snake is paddle-shaped and adapted to swimming.
On average, the total length of a male is long, with a long tail. Females are significantly larger, with an average total length of and a tail length of .
Distribution and habitat
The yellow-lipped sea krait is widespread throughout the eastern Indian Ocean and Western Pacific. It can be found from the eastern coast of India, along the coast of the Bay of Bengal in Bangladesh, Myanmar, and other parts of Southeast Asia, to the Malay Archipelago and to some parts of southern China, Taiwan, and the Ryukyu Islands of Japan. The species is also common near Fiji and other Pacific islands within its range. Vagrant individuals have been recorded in Australia, New Caledonia, and New Zealand. Six specimens have been found around the North Island of New Zealand between 1880 and 2005, suspected to have come from populations based in Fiji and Tonga. It is the most common sea krait identified in New Zealand, and second-most seen sea snake after the yellow-bellied sea snake - common enough to be considered a native species, protected under the Wildlife Act 1953.
Venom
The venom of this elapid, L. colubrina, is a very powerful neurotoxic protein, with a subcutaneous LD50 in mice of 0.45 mg/kg body weight. The venom is an α-neurotoxin that disrupts synapses by competing with acetylcholine for receptors on the postsynaptic membrane, similar to erabutoxins and α-bungarotoxins. In mice, lethal venom doses cause lethargy, flaccid paralysis, and convulsions in quick succession before death. Dogs injected with lethal doses produced symptoms consistent with fatal hypertension and cyanosis observed in human sea snake bite victims.
Some varieties of eels, which are a primary food source for yellow-lipped sea kraits, may have coevolved resistance to yellow-lipped sea krait venom. Gymnothorax moray eels taken from the Caribbean, where yellow-lipped sea kraits are not endemic, died after injection with doses as small as 0.1 mg/kg body weight, but Gymnothorax individuals taken from New Guinea, where yellow-lipped sea kraits are endemic, were able to tolerate doses as large as 75 mg/kg without severe injury.
Behavior
Yellow-lipped sea kraits are semiaquatic. Juveniles stay in water and on adjacent coasts, but adults are able to move further inland and spend half their time on land and half in the ocean. Adult males are more terrestrially active during mating and hunt in shallower water, requiring more terrestrial locomotive ability. Adult females, though, are less active on land during mating and hunt in deeper water, requiring more aquatic locomotive ability. Because males are smaller, they crawl and swim faster than females.
Body adaptations, especially a paddle-like tail, help yellow-lipped sea kraits to swim. These adaptations are also found in more distantly related sea snakes (Hydrophiinae) because of convergent evolution, but because of the differences in motion between crawling and swimming, these same adaptations impede the snake's terrestrial motion. On dry land, a yellow-lipped sea krait can still move, but typically at only slightly more than a fifth of its swimming speed. In contrast, most sea snakes other than Laticauda spp. are virtually stranded on dry land.
When hunting, yellow-lipped sea kraits frequently head into deep water far from land, but return to land to digest meals, shed skin, and reproduce. Individuals return to their specific home islands, exhibiting philopatry. When yellow-lipped sea kraits on Fijian islands were relocated to different islands 5.3 km away, all recaptured individuals were found on their home islands in an average of 30.7 days.
Yellow-lipped sea kraits collected near the tip of Borneo had heavy tick infections.
Hunting and diet
Hunting is often performed alone, but L. colubrina kraits may also do so in large numbers in the company of hunting parties of giant trevally and goatfish. This cooperative hunting technique is similar to that of the moray eel, with the yellow-lipped sea kraits flushing out prey from narrow crevices and holes, and the trevally and goatfish feeding on fleeing prey.
While probing crevices with their heads, yellow-lipped sea kraits are unable to observe approaching predators and can be vulnerable. The snakes can deter predators, such as larger fish, sharks, and birds, by fooling them into thinking that their tail is their head, because the color and movement of the tail is similar to that of the snake's head. For example, the lateral aspect of the tail corresponds to the dorsal view of the head.
Yellow-lipped sea kraits primarily feed on varieties of eels (of the families Congridae, Muraenidae, and Ophichthidae), but also eat small fish (including those of the families Pomacentridae and Synodontidae). Males and females exhibit sexual dimorphism in hunting behavior, as adult females, which are significantly larger than males, prefer to hunt in deeper water for larger conger eels, while adult males hunt in shallower water for smaller moray eels. In addition, females hunt for only one prey item per foraging bout, while males often hunt for multiple items. After hunting, yellow-lipped sea kraits return to land to digest their prey.
Courtship and reproduction
The yellow-lipped sea krait is oviparous, meaning it lays eggs that develop outside of the body.
Each year during the warmer months of September through December, males gather on land and in the water around gently sloping areas at high tide. Males prefer to mate with larger females because they produce larger and more offspring.
When a male detects a female, he chases the female and begins courtship. Females are larger and slower than males, and many males escort and intertwine around a single female. The males then align their bodies with the female and rhythmically contract; the resulting mass of snakes can remain nearly motionless for several days. After courtship, the snakes copulate for about an average of two hours.
The female yellow-lipped sea kraits then lay as many as 10 eggs per clutch. The eggs are deposited in crevices where they remain until hatching. These eggs are very rarely found in the wild; only two nests have been definitively reported throughout the entire range of the species.
Interaction with humans
Because yellow-lipped sea kraits spend much of their time on land, they are often encountered by humans. They are frequently found in the water intake and exhaust pipes of boats. They are also attracted to light and can be distracted by artificial sources of light, including hotels and other buildings, on coasts.
Fewer bites from this species are recorded compared to other venomous species such as cobras and vipers, as it is less aggressive and tends to avoid humans. If they do bite, it is usually in self-defense when accidentally grabbed. Most sea snake bites occur when fishermen attempt to untangle the snakes from their fishing nets.
In the Philippines, yellow-lipped sea kraits are caught for their skin and meat; the meat is smoked and exported for use in Japanese cuisine. The smoked meat of a related Laticauda species, the black-banded sea krait, is used in Okinawan cuisine to make irabu-jiru (, irabu soup).
| Biology and health sciences | Snakes | Animals |
3739933 | https://en.wikipedia.org/wiki/Turing%27s%20proof | Turing's proof | Turing's proof is a proof by Alan Turing, first published in November 1936 with the title "On Computable Numbers, with an Application to the ". It was the second proof (after Church's theorem) of the negation of Hilbert's ; that is, the conjecture that some purely mathematical yes–no questions can never be answered by computation; more technically, that some decision problems are "undecidable" in the sense that there is no single algorithm that infallibly gives a correct "yes" or "no" answer to each instance of the problem. In Turing's own words:
"what I shall prove is quite different from the well-known results of Gödel ... I shall now show that there is no general method which tells whether a given formula U is provable in K [Principia Mathematica]".
Turing followed this proof with two others. The second and third both rely on the first. All rely on his development of typewriter-like "computing machines" that obey a simple set of rules and his subsequent development of a "universal computing machine". As per UK copyright law, the novella entered the public domain on 1 January 2025, 70 full calendar years after Turing's death on 7 June 1954.
Summary of the proofs
In his proof that the Entscheidungsproblem can have no solution, Turing proceeded from two proofs that were to lead to his final proof. His first theorem is most relevant to the halting problem, the second is more relevant to Rice's theorem.
First proof: that no "computing machine" exists that can decide whether or not an arbitrary "computing machine" (as represented by an integer 1, 2, 3, . . .) is "circle-free" (i.e. goes on printing its number in binary ad infinitum): "...we have no general process for doing this in a finite number of steps" (p. 132, ibid.). Turing's proof, although it seems to use the "diagonal process", in fact shows that his machine (called H) cannot calculate its own number, let alone the entire diagonal number (Cantor's diagonal argument): "The fallacy in the argument lies in the assumption that B [the diagonal number] is computable" The proof does not require much mathematics.
Second proof: This one is perhaps more familiar to readers as Rice's theorem: "We can show further that there can be no machine E which, when supplied with the S.D ["program"] of an arbitrary machine M, will determine whether M ever prints a given symbol (0 say)"
Third proof: "Corresponding to each computing machine M we construct a formula Un(M) and we show that, if there is a general method for determining whether Un(M) is provable, then there is a general method for determining whether M ever prints 0".
The third proof requires the use of formal logic to prove a first lemma, followed by a brief word-proof of the second:
Finally, in only 64 words and symbols Turing proves by reductio ad absurdum that "the Hilbert Entscheidungsproblem can have no solution".
Summary of the first proof
Turing created a thicket of abbreviations. See the glossary at the end of the article for definitions.
Some key clarifications:
Turing spent much of his paper actually "constructing" his machines to convince us of their truth. This was required by his use of the reductio ad absurdum form of proof. We must emphasize the "constructive" nature of this proof. Turing describes what could be a real machine, really buildable. The only questionable element is the existence of machine D, which this proof will eventually show to be impossible.
Turing begins the proof with the assertion of the existence of a “decision/determination” machine D. When fed any S.D (string of symbols A, C, D, L, R, N, semicolon “;”) it will determine if this S.D (symbol string) represents a "computing machine" that is either "circular" — and therefore "un-satisfactory u" — or "circle-free" — and therefore "satisfactory s".
Turing makes no comment about how machine D goes about its work. For sake of argument, we suppose that D would first look to see if the string of symbols is "well-formed" (i.e. in the form of an algorithm and not just a scramble of symbols), and if not then discard it. Then it would go “circle-hunting”. To do this perhaps it would use “heuristics” (tricks: taught or learned). For purposes of the proof, these details are not important.
Turing then describes (rather loosely) the algorithm (method) to be followed by a machine he calls H. Machine H contains within it the decision-machine D (thus D is a “subroutine” of H). Machine H’s algorithm is expressed in H’s table of instructions, or perhaps in H’s Standard Description on tape and united with the universal machine U; Turing does not specify this.
Machine H is responsible for converting any number N into an equivalent S.D symbol string for sub-machine D to test. (In programming parlance: H passes an arbitrary "S.D” to D, and D returns “satisfactory” or “unsatisfactory”.) Machine H is also responsible for keeping a tally R (“Record”?) of successful numbers (we suppose that the number of “successful” S.D's, i.e. R, is much less than the number of S.D's tested, i.e. N). Finally, H prints on a section of its tape a diagonal number “beta-primed” B’. H creates this B’ by “simulating” (in the computer-sense) the “motions” of each “satisfactory” machine/number; eventually this machine/number under test will arrive at its Rth “figure” (1 or 0), and H will print it. H then is responsible for “cleaning up the mess” left by the simulation, incrementing N and proceeding onward with its tests, ad infinitum.
Note: All these machines that H is hunting for are what Turing called "computing machines". These compute binary-decimal-numbers in an endless stream of what Turing called "figures": only the symbols 1 and 0.
An example to illustrate the first proof
An example: Suppose machine H has tested 13472 numbers and produced 5 satisfactory numbers, i.e. H has converted the numbers 1 through 13472 into S.D's (symbol strings) and passed them to D for test. As a consequence H has tallied 5 satisfactory numbers and run the first one to its 1st "figure", the second to its 2nd figure, the third to its 3rd figure, the fourth to its 4th figure, and the fifth to its 5th figure. The count now stands at N = 13472, R = 5, and B' = ".10011" (for example). H cleans up the mess on its tape, and proceeds:
H increments N = 13473 and converts "13473" to symbol string ADRLD. If sub-machine D deems ADLRD unsatisfactory, then H leaves the tally-record R at 5. H will increment the number N to 13474 and proceed onward. On the other hand, if D deems ADRLD satisfactory then H will increment R to 6. H will convert N (again) into ADLRD [this is just an example, ADLRD is probably useless] and “run” it using the universal machine U until this machine-under-test (U "running" ADRLD) prints its 6th “figure” i.e. 1 or 0. H will print this 6th number (e.g. “0”) in the “output” region of its tape (e.g. B’ = “.100110”).
H cleans up the mess, and then increments the number N to 13474.
The whole process unravels when H arrives at its own number K. We will proceed with our example. Suppose the successful-tally/record R stands at 12. H finally arrives at its own number minus 1, i.e. N = K-1 = 4335...3214, and this number is unsuccessful. Then H increments N to produce K = 4355...3215, i.e. its own number. H converts this to “LDDR...DCAR” and passes it to decision-machine D. Decision-machine D must return “satisfactory” (that is: H must by definition go on and on testing, ad infinitum, because it is "circle-free"). So H now increments tally R from 12 to 13 and then re-converts the number-under-test K into its S.D and uses U to simulate it. But this means that H will be simulating its own motions. What is the first thing the simulation will do? This simulation K-aka-H either creates a new N or “resets” the “old” N to 1. This "K-aka-H" either creates a new R or “resets” the “old” R to 0. Old-H “runs” new "K-aka-H" until it arrives at its 12th figure.
But it never makes it to the 13th figure; K-aka-H eventually arrives at 4355...3215, again, and K-aka-H must repeat the test. K-aka-H will never reach the 13th figure. The H-machine probably just prints copies of itself ad infinitum across blank tape. But this contradicts the premise that H is a satisfactory, non-circular computing machine that goes on printing the diagonal numbers's 1's and 0's forever. (We will see the same thing if N is reset to 1 and R is reset to 0.)
If the reader does not believe this, they can write a "stub" for decision-machine D (stub "D" will return "satisfactory") and then see for themselves what happens at the instant machine H encounters its own number.
Summary of the second proof
Less than one page long, the passage from premises to conclusion is obscure.
Turing proceeds by reductio ad absurdum. He asserts the existence of a machine E, which when given the S.D (Standard Description, i.e. "program") of an arbitrary machine M, will determine whether M ever prints a given symbol (0 say). He does not assert that this M is a "computing machine".
Given the existence of machine E, Turing proceeds as follows:
If machine E exists then a machine G exists that determines if M prints 0 infinitely often, AND
If E exists then another process exists [we can call the process/machine G' for reference] that determines if M prints 1 infinitely often, THEREFORE
When we combine G with G' we have a process that determines if M prints an infinity of figures, AND
IF the process "G with G'" determines M prints an infinity of figures, THEN "G with G'" has determined that M is circle-free, BUT
This process "G with G'" that determine if M is circle-free, by proof 1, cannot exist, THEREFORE
Machine E does not exist.
Details of second proof
The difficulty in the proof is step 1. The reader will be helped by realizing that Turing is not explaining his subtle handiwork. (In a nutshell: he is using certain equivalencies between the “existential-“ and “universal-operators” together with their equivalent expressions written with logical operators.)
Here's an example: Suppose we see before us a parking lot full of hundreds of cars. We decide to go around the entire lot looking for: “Cars with flat (bad) tires”. After an hour or so we have found two “cars with bad tires.” We can now say with certainty that “Some cars have bad tires”. Or we could say: “It’s not true that ‘All the cars have good tires’”. Or: “It is true that: ‘not all the cars have good tires”. Let us go to another lot. Here we discover that “All the cars have good tires.” We might say, “There’s not a single instance of a car having a bad tire.” Thus we see that, if we can say something about each car separately then we can say something about ALL of them collectively.
This is what Turing does:
From M he creates a collection of machines {M1, M2, M3, M4, ..., Mn} and about each he writes a sentence: “X prints at least one 0” and allows only two “truth values”, True = blank or False = :0:. One by one he determines the truth value of the sentence for each machine and makes a string of blanks or :0:, or some combination of these. We might get something like this: “M1 prints a 0” = True AND “M2 prints a 0” = True AND “M3 prints a 0” = True AND “M4 prints a 0” = False, ... AND “Mn prints a 0” = False. He gets the string
if there are an infinite number of machines Mn. If on the other hand if every machine had produced a "True" then the expression on the tape would be
Thus Turing has converted statements about each machine considered separately into a single "statement" (string) about all of them. Given the machine (he calls it G) that created this expression, he can test it with his machine E and determine if it ever produces a 0. In our first example above we see that indeed it does, so we know that not all the M's in our sequence print 0s. But the second example shows that, since the string is blanks then every Mn in our sequence has produced a 0.
All that remains for Turing to do is create a process to create the sequence of Mn's from a single M.
Suppose M prints this pattern:
M => ...AB01AB0010AB…
Turing creates another machine F that takes M and crunches out a sequence of Mn's that successively convert the first n 0's to “0-bar” (0):
He states, without showing details, that this machine F is truly build-able. We can see that one of a couple things could happen. F may run out of machines that have 0's, or it may have to go on ad infinitum creating machines to “cancel the zeros”.
Turing now combines machines E and F into a composite machine G. G starts with the original M, then uses F to create all the successor-machines M1, M2,. . ., Mn. Then G uses E to test each machine starting with M. If E detects that a machine never prints a zero, G prints :0: for that machine. If E detects that a machine does print a 0 (we assume, Turing doesn’t say) then G prints :: or just skips this entry, leaving the squares blank. We can see that a couple things can happen.
Now, what happens when we apply E to G itself?
As we can apply the same process for determining if M prints 1 infinitely often. When we combine these processes, we can determine that M does, or does not, go on printing 1's and 0's ad infinitum. Thus we have a method for determining if M is circle-free. By Proof 1 this is impossible. So the first assertion that E exists, is wrong: E does not exist.
Summary of the third proof
Here Turing proves "that the Hilbert Entscheidungsproblem can have no solution". Here he
Both Lemmas #1 and #2 are required to form the necessary "IF AND ONLY IF" (i.e. logical equivalence) required by the proof:
Turing demonstrates the existence of a formula Un(M) which says, in effect, that "in some complete configuration of M, 0 appears on the tape" (p. 146). This formula is TRUE, that is, it is "constructible", and he shows how to go about this.
Then Turing proves two Lemmas, the first requiring all the hard work. (The second is the converse of the first.) Then he uses reductio ad absurdum to prove his final result:
There exists a formula Un(M). This formula is TRUE, AND
If the Entscheidungsproblem can be solved THEN a mechanical process exists for determining whether Un(M) is provable (derivable), AND
By Lemmas 1 and 2: Un(M) is provable IF AND ONLY IF 0 appears in some "complete configuration" of M, AND
IF 0 appears in some "complete configuration" of M THEN a mechanical process exists that will determine whether arbitrary M ever prints 0, AND
By Proof 2 no mechanical process exists that will determine whether arbitrary M ever prints 0, THEREFORE
Un(M) is not provable (it is TRUE, but not provable) which means that the Entscheidungsproblem is unsolvable.
Details of the third proof
[If readers intend to study the proof in detail they should correct their copies of the pages of the third proof with the corrections that Turing supplied. Readers should also come equipped with a solid background in (i) logic (ii) the paper of Kurt Gödel: "On Formally Undecidable Propositions of Principia Mathematica and Related Systems". For assistance with Gödel's paper they may consult e.g. Ernest Nagel and James R. Newman, Gödel's Proof, New York University Press, 1958.]
To follow the technical details, the reader will need to understand the definition of "provable" and be aware of important "clues".
"Provable" means, in the sense of Gödel, that (i) the axiom system itself is powerful enough to produce (express) the sentence "This sentence is provable", and (ii) that in any arbitrary "well-formed" proof the symbols lead by axioms, definitions, and substitution to the symbols of the conclusion.
First clue: "Let us put the description of M into the first standard form of §6". Section 6 describes the very specific "encoding" of machine M on the tape of a "universal machine" U. This requires the reader to know some idiosyncrasies of Turing's universal machine U and the encoding scheme.
(i) The universal machine is a set of "universal" instructions that reside in an "instruction table". Separate from this, on U's tape, a "computing machine" M will reside as "M-code". The universal table of instructions can print on the tape the symbols A, C, D, 0, 1, u, v, w, x, y, z, : . The various machines M can print these symbols only indirectly by commanding U to print them.
(ii) The "machine code" of M consists of only a few letters and the semicolon, i.e. D, C, A, R, L, N, ; . Nowhere within the "code" of M will the numerical "figures" (symbols) 1 and 0 ever appear. If M wants U to print a symbol from the collection blank, 0, 1 then it uses one of the following codes to tell U to print them. To make things more confusing, Turing calls these symbols S0, S1, and S2, i.e.
blank = S0 = D
0 = S1 = DC
1 = S2 = DCC
(iii) A "computing machine", whether it is built directly into a table (as his first examples show), or as machine-code M on universal-machine U's tape, prints its number on blank tape (to the right of M-code, if there is one) as 1s and 0s forever proceeding to the right.
(iv) If a "computing machine" is U+"M-code", then "M-code" appears first on the tape; the tape has a left end and the "M-code" starts there and proceeds to the right on alternate squares. When the M-code comes to an end (and it must, because of the assumption that these M-codes are finite algorithms), the "figures" will begin as 1s and 0s on alternate squares, proceeding to the right forever. Turing uses the (blank) alternate squares (called "E"- "eraseable"- squares) to help U+"M-code" keep track of where the calculations are, both in the M-code and in the "figures" that the machine is printing.
(v) A "complete configuration" is a printing of all symbols on the tape, including M-code and "figures" up to that point, together with the figure currently being scanned (with a pointer-character printed to the left of the scanned symbol?). If we have interpreted Turing's meaning correctly, this will be a hugely long set of symbols. But whether the entire M-code must be repeated is unclear; only a printing of the current M-code instruction is necessary plus the printing of all figures with a figure-marker).
(vi) Turing reduced the vast possible number of instructions in "M-code" (again: the code of M to appear on the tape) to a small canonical set, one of three similar to this: {qi Sj Sk R ql} e.g. If machine is executing instruction #qi and symbol Sj is on the square being scanned, then Print symbol Sk and go Right and then go to instruction ql: The other instructions are similar, encoding for "Left" L and "No motion" N. It is this set that is encoded by the string of symbols qi = DA...A, Sj = DC...C, Sk = DC...C, R, ql = DA....A. Each instruction is separated from another one by the semicolon. For example, {q5, S1 S0 L q3} means: Instruction #5: If scanned symbol is 0 then print blank, go Left, then go to instruction #3. It is encoded as follows
Second clue: Turing is using ideas introduced in Gödel's paper, that is, the "Gödelization" of (at least part of) the formula for Un(M). This clue appears only as a footnote on page 138 (): "A sequence of r primes is denoted by ^(r)" (ibid.) [Here, r inside parentheses is "raised".] This "sequence of primes" appears in a formula called F^(n).
Third clue: This reinforces the second clue. Turing's original attempt at the proof uses the expression:
Earlier in the paper Turing had previously used this expression (p. 138) and defined N(u) to mean "u is a non-negative integer" (ibid.) (i.e. a Gödel number). But, with the Bernays corrections, Turing abandoned this approach (i.e. the use of N(u)) and the only place where "the Gödel number" appears explicitly is where he uses F^(n).
What does this mean for the proof? The first clue means that a simple examination of the M-code on the tape will not reveal if a symbol 0 is ever printed by U+"M-code". A testing-machine might look for the appearance of DC in one of the strings of symbols that represent an instruction. But will this instruction ever be "executed?" Something has to "run the code" to find out. This something can be a machine, or it can be lines in a formal proof, i.e. Lemma #1.
The second and third clues mean that, as its foundation is Gödel's paper, the proof is difficult.
In the example below we will actually construct a simple "theorem"—a little Post–Turing machine program "run it". We will see just how mechanical a properly designed theorem can be. A proof, we will see, is just that, a "test" of the theorem that we do by inserting a "proof example" into the beginning and see what pops out at the end.
Both Lemmas #1 and #2 are required to form the necessary "IF AND ONLY IF" (i.e. logical equivalence) required by the proof:
To quote Franzén:
Franzén has defined "provable" earlier in his book:
Thus a "sentence" is a string of symbols, and a theorem is a string of strings of symbols.
Turing is confronted with the following task:
Thus the "string of sentences" will be strings of strings of symbols. The only allowed individual symbols will come from Gödel's symbols defined in his paper.(In the following example we use the "<" and ">" around a "figure" to indicate that the "figure" is the symbol being scanned by the machine).
An example to illustrate the third proof
In the following, we have to remind ourselves that every one of Turing's “computing machines” is a binary-number generator/creator that begins work on “blank tape”. Properly constructed, it always cranks away ad infinitum, but its instructions are always finite. In Turing's proofs, Turing's tape had a “left end” but extended right ad infinitum. For sake of example below we will assume that the “machine” is not a Universal machine, but rather the simpler “dedicated machine” with the instructions in the Table.
Our example is based on a modified Post–Turing machine model of a Turing Machine. This model prints only the symbols 0 and 1. The blank tape is considered to be all b's. Our modified model requires us to add two more instructions to the 7 Post–Turing instructions. The abbreviations that we will use are:
In the cases of R, L, E, P0, and P1 after doing its task the machine continues on to the next instruction in numerical sequence; ditto for the jumps if their tests fail.
But, for brevity, our examples will only use three squares. And these will always start as there blanks with the scanned square on the left: i.e. bbb. With two symbols 1, 0 and blank we can have 27 distinct configurations:
We must be careful here, because it is quite possible that an algorithm will (temporarily) leave blanks in between figures, then come back and fill something in. More likely, an algorithm may do this intentionally. In fact, Turing's machine does this—it prints on alternate squares, leaving blanks between figures so it can print locator symbols.
Turing always left alternate squares blank so his machine could place a symbol to the left of a figure (or a letter if the machine is the universal machine and the scanned square is actually in the “program”). In our little example we will forego this and just put symbols ( ) around the scanned symbol, as follows:
Let us write a simple program:
Remember that we always start with blank tape. The complete configuration prints the symbols on the tape followed by the next instruction:
Let us add “jump” into the formula. When we do this we discover why the complete configuration must include the tape symbols. (Actually, we see this better, below.) This little program prints three “1”s to the right, reverses direction and moves left printing 0’s until it hits a blank. We will print all the symbols that our machine uses:
Here at the end we find that a blank on the left has “come into play” so we leave it as part of the total configuration.
Given that we have done our job correctly, we add the starting conditions and see “where the theorem goes”. The resulting configuration—the number 110—is the PROOF.
Turing's first task had to write a generalized expression using logic symbols to express exactly what his Un(M) would do.
Turing's second task is to "Gödelize" this hugely long string-of-string-of-symbols using Gödel's technique of assigning primes to the symbols and raising the primes to prime-powers, per Gödel's method.
Complications
Turing's proof is complicated by a large number of definitions, and confounded with what Martin Davis called "petty technical details" and "...technical details [that] are incorrect as given". Turing himself published "A Correction" in 1938: "The author is indebted to P. Bernays for pointing out these errors".
Specifically, in its original form the third proof is badly marred by technical errors. And even after Bernays' suggestions and Turing's corrections, errors remained in the description of the universal machine. And confusingly, since Turing was unable to correct his original paper, some text within the body harks to Turing's flawed first effort.
Bernays' corrections may be found in ; the original is to be found as "On Computable Numbers, with an Application to the Entscheidungsproblem. A Correction," Proceedings of the London Mathematical Society (2), 43 (1938), 544-546.
The on-line version of Turing's paper has these corrections in an addendum; however, corrections to the Universal Machine must be found in an analysis provided by Emil Post.
At first, the only mathematician to pay close attention to the details of the proof was Post (cf. Hodges p. 125) — mainly because he had arrived simultaneously at a similar reduction of "algorithm" to primitive machine-like actions, so he took a personal interest in the proof. Strangely (perhaps World War II intervened) it took Post some ten years to dissect it in the Appendix to his paper Recursive Unsolvability of a Problem of Thue, 1947.
Other problems present themselves: In his Appendix Post commented indirectly on the paper's difficulty and directly on its "outline nature" and "intuitive form" of the proofs. Post had to infer various points:
Anyone who has ever tried to read the paper will understand Hodges' complaint:
Glossary of terms used by Turing
1 computable number — a number whose decimal is computable by a machine (i.e., by finite means such as an algorithm)
2 M — a machine with a finite instruction table and a scanning/printing head. M moves an infinite tape divided into squares each “capable of bearing a symbol”. The machine-instructions are only the following: move one square left, move one square right, on the scanned square print symbol p, erase the scanned square, if the symbol is p then do instruction aaa, if the scanned symbol is not p then do instruction aaa, if the scanned symbol is none then do instruction aaa, if the scanned symbol is any do instruction aaa [where “aaa” is an instruction-identifier].
3 computing machine — an M that prints two kinds of symbols, symbols of the first type are called “figures” and are only binary symbols 1 and 0; symbols of the second type are any other symbols.
4 figures — symbols 1 and 0, a.k.a. “symbols of the first kind”
5 m-configuration — the instruction-identifier, either a symbol in the instruction table, or a string of symbols representing the instruction- number on the tape of the universal machine (e.g. "DAAAAA = #5")
6 symbols of the second kind — any symbols other than 1 and 0
7 circular — an unsuccessful computating machine. It fails to print, ad infinitum, the figures 0 or 1 that represent in binary the number it computes
8 circle-free — a successful computating machine. It prints, ad infinitum, the figures 0 or 1 that represent in binary the number it computes
9 sequence — as in “sequence computed by the machine”: symbols of the first kind a.k.a. figures a.k.a. symbols 0 and 1.
10 computable sequence — can be computed by a circle-free machine
11 S.D – Standard Description: a sequence of symbols A, C, D, L, R, N, “;” on a Turing machine tape
12 D.N — Description number: an S.D converted to a number: 1=A, 2=C, 3 =D, 4=L, 5=R, 6=N, 7=;
13 M(n) — a machine whose D.N is number “n”
14 satisfactory — a S.D or D.N that represents a circle-free machine
15 U — a machine equipped with a “universal” table of instructions. If U is “supplied with a tape on the beginning of which is written the S.D of some computing machine M, U will compute the same sequence as M.”
16 β’—“beta-primed”: A so-called “diagonal number” made up of the n-th figure (i.e. 0 or 1) of the n-th computable sequence [also: the computable number of H, see below]
17 u — an unsatisfactory, i.e. circular, S.D
18 s — satisfactory, i.e. circle-free S.D
19 D — a machine contained in H (see below). When supplied with the S.D of any computing machine M, D will test M's S.D and if circular mark it with “u” and if circle-free mark it with “s”
20 H — a computing machine. H computes B’, maintains R and N. H contains D and U and an unspecified machine (or process) that maintains N and R and provides D with the equivalent S.D of N. E also computes the figures of B’ and assembles the figures of B’.
21 R — a record, or tally, of the quantity of successful (circle-free) S.D tested by D
22 N — a number, starting with 1, to be converted into an S.D by machine E. E maintains N.
23 K — a number. The D.N of H.
Required for Proof #3
5 m-configuration — the instruction-identifier, either a symbol in the instruction table, or a string of symbols representing the instruction's number on the tape of the universal machine (e.g. "DAAAAA = instruction #5"). In Turing's S.D the m-configuration appears twice in each instruction, the left-most string is the "current instruction"; the right-most string is the next instruction.
24 complete configuration — the number (figure 1 or 0) of the scanned square, the complete sequence of all symbols on the tape, and the m-configuration (the instruction-identifier, either a symbol or a string of symbols representing a number, e.g. "instruction DAAAA = #5")
25 RSi(x, y) — "in the complete configuration x of M the symbol on square y is Si; "complete configuration" is definition #5
26 I(x, y) — "in the complete configuration x of M the square y is scanned"
27 Kqm(x) — "in the complete configuration x of M the machine-configuration (instruction number) is qm"
28 F(x,y) — "y is the immediate successor of x" (follows Gödel's use of "f" as the successor-function).
29 G(x,y) — "x precedes y", not necessarily immediately
30 Inst{qi, Sj Sk L ql} is an abbreviation, as are Inst{qi, Sj Sk R ql}, and Inst{qi, Sj Sk N ql}. See below.
Turing reduces his instruction set to three “canonical forms” – one for Left, Right, and No-movement. Si and Sk are symbols on the tape.
For example, the operations in the first line are PSk = PRINT symbol Sk from the collection A, C, D, 0, 1, u, v, w, x, y, z, :, then move tape LEFT.
These he further abbreviated as:
(N1) qi Sj Sk L qm
(N2) qi Sj Sk R qm
(N3) qi Sj Sk N qm
In Proof #3 he calls the first of these “Inst{qi Sj Sk L ql}”, and he shows how to write the entire machine S.D as the logical conjunction (logical OR): this string is called “Des(M)”, as in “Description-of-M”.
i.e. if the machine prints 0 then 1's and 0's on alternate squares to the right ad infinitum it might have the table (a similar example appears on page 119):
(This has been reduced to canonical form with the “p-blank” instructions so it differs a bit from Turing's example.)
If put them into the “ Inst( ) form” the instructions will be the following (remembering: S0 is blank, S1 = 0, S2 = 1):
The reduction to the Standard Description (S.D) will be:
This agrees with his example in the book (there will be a blank between each letter and number). Universal machine U uses the alternate blank squares as places to put "pointers".
| Mathematics | Computability theory | null |
3741200 | https://en.wikipedia.org/wiki/Asian%20water%20monitor | Asian water monitor | The Asian water monitor (Varanus salvator) is a large varanid lizard native to South and Southeast Asia. It is widely considered to be the second-largest lizard species, after the Komodo dragon. It is distributed from eastern and northeastern India and Bangladesh, the Andaman and Nicobar Islands, Sri Lanka, through southern China and Hainan Island in the east to mainland Southeast Asia and the islands of Sumatra, Borneo, Java, Lombok, the Riau Archipelago, Sulawesi. It is one of the most widespread monitor lizards.
The Asian water monitor has a natural affinity towards water, inhabiting the surroundings of lakes, rivers, ponds, swamps and various riparian habitats, including sewers, city parks, and urban waterways. It is an excellent swimmer and hunts fish, frogs, invertebrates, water birds, and other types of aquatic and amphibious prey.
Due to its apparently large, stable population, it is currently listed as Least Concern on to the IUCN Red List.
Etymology
The generic name Varanus is derived from the Arabic waral (), which translates as "monitor". The specific name is the Latin word for "saviour", denoting a possible religious connotation. The water monitor is occasionally confused with the crocodile monitor (V. salvadorii) because of their similar scientific names.
Some common names for the species are Malayan water monitor, common water monitor, two-banded monitor, rice lizard, ring lizard, plain lizard, no-mark lizard and water monitor etc.
Taxonomy
Stellio salvator was the scientific name used by Josephus Nicolaus Laurenti in 1768 when he described a water monitor specimen.
The family Varanidae contains nearly 80 species of monitor lizards, all of which belong to the genus Varanus. There is a significant amount of taxonomic uncertainty within this species complex. Morphological analyses have begun to unravel this taxonomic uncertainty but molecular studies are needed to test and confirm the validity of certain groupings within this genus. Research initiatives such as these are very important to assess changes in conservation assessments.
Subspecies
V. s. salvator is the nominate subspecies and is restricted to Sri Lanka.
V. s. andamanensis, the Andaman Islands water monitor, inhabits the Andaman Islands and the southern Nicobar Islands; the type locality is Port Blair.
V. s. bivittatus (Mertens 1959), the two-striped water monitor, is common to the Indonesian islands of Java, Bali, Lombok, Sumbawa, Flores, Alor, Wetar and some neighbouring islands within the Sunda archipelago. The type locality is Java.
V. s. macromaculatus, the Southeast Asian water monitor, occurs from Bihar, West Bengal to Assam, Arunachal Pradesh, Manipur and Meghalaya through Bangladesh, Myanmar, Thailand, Cambodia, Laos, Vietnam to Malaysia, Singapore, Sumatra, Borneo, Bali, Bangka, Batam, Belitung, Bintan, Enggano, Matak, Nikoi. The type specimen was captured in Thailand.
V. s. ziegleri, Ziegler's water monitor, occurs on the Obi Islands, North Maluku, Indonesia.
V. s. celebensis, the Sulawesi water monitor, occurs primarily in North Sulawesi Province of the island of Sulawesi, Indonesia.
Varanus cumingi, Varanus marmoratus, and Varanus nuchalis were classified as subspecies until 2007, when they were elevated to full species.
The black water monitor from Thailand's Satun Province and Thai-Malaysian border area was formerly the subspecies V. s. komaini, but now is regarded as a junior synonym and melanistic population of V. s. macromaculatus.
Description
The Asian water monitor is dark brown or blackish with yellow spots on the underside that fade gradually with age. It has blackish bands with yellow edges extending back from each eye. Its body is muscular, with long, powerful, laterally compressed tails. Its scales are keeled; the ones on top of the head are larger than those on the back. Its neck is long and the snout elongated. It has powerful jaws, serrated teeth and sharp claws.
Adults rarely exceed in length, but the largest specimen on record from Sri Lanka measured . A common mature weight is . However, 80 males killed for the leather trade in Sumatra averaged only and snout-to-vent and in total length; 42 females averaged and snout–vent length and in total length.
Males are larger than females and attain breeding maturity at a length of and a weight of ; and females at a length of .
A series of adults weighed .
Mature individuals in northern Sumatra were estimated to have a mean estimated body mass of .
A sample of 55 Asian water monitors weighed . The maximum weight of captive individuals is over .
In captivity, Asian water monitors' life expectancy has been determined to be anywhere between 11 and 25 years depending on conditions, in the wild it is considerably shorter.
The teeth are compressed, serrated (though irregularly) and recurved. Up to two replacement teeth lie behind each tooth position at a given time, and teeth are replaced every 59 days.
Distribution and habitat
The Asian water monitor is widely distributed from India, Bangladesh, Sri Lanka, Myanmar and Thailand, Cambodia, Laos, Vietnam, the Chinese Guangxi and Hainan provinces, Malaysia, Singapore to the Sunda islands Sumatra, Java, Bali, Borneo and Sulawesi. It inhabits primarily lowland freshwater and brackish wetlands. It has been recorded up to an elevation of .
The Asian water monitor is semiaquatic and opportunistic; it inhabits a variety of natural habitats though predominantly resides in primary forests and mangrove swamps. It has been noted that it is not deterred from living in areas near human civilization. In fact, it has been known to adapt and thrive in agricultural areas as well as cities with canal systems, such as in Sri Lanka, where they are not hunted or persecuted. Habitats that are considered to be most important are mangrove vegetation, swamps, wetlands, and elevations below . It does not thrive in habitats with extensive loss of natural vegetation and aquatic resources.
A population of these monitors have become established as an invasive population in the southeastern parts of the USA.
Behaviour and ecology
Water monitors defend themselves using their tails, claws, and jaws. They are excellent swimmers, using the raised fin on their tails to steer through water. When encountering smaller prey items, the water monitor will subdue it in its jaws and proceed to violently thrash its neck, destroying the prey's organs and spine which leaves it dead or incapacitated. The lizard will then swallow it whole.
In dominantly aquatic habitats, their semiaquatic behavior is considered to provide a measure of safety from predators. This along with their versatile diet is said to contribute to their plasticity, or ecological adaptability. When hunted by predators such as the king cobra (Ophiophagus hannah) they will climb trees using their powerful legs and claws. If this evasion is not enough to escape danger, they have also been known to jump from trees into streams for safety, a tactic similar to that of the green iguana (Iguana iguana).
Like the Komodo dragon, the water monitor will often eat carrion, or rotten flesh. By eating this decaying flesh, the lizard provides benefits to the ecosystem by removing infectious elements, cleaning the environment. They have a keen sense of smell and can smell a carcass from far away.
While adults are terrestrial, juveniles are primarily arboreal.
The first description of the water monitor and its behaviour in English literature was made in 1681 by Robert Knox, who observed it during his long confinement in the Kingdom of Kandy: "There is a Creature here called Kobberaguion, resembling an Alligator. The biggest may be five or six feet long, speckled black and white. He lives most upon the Land, but will take the water and dive under it: hath a long blue forked tongue like a sting, which he puts forth and hisseth and gapeth, but doth not bite nor sting, tho the appearance of him would scare those that knew not what he was. He is not afraid of people, but will lie gaping and hissing at them in the way, and will scarce stir out of it. He will come and eat Carrion with the Dogs and Jackals, and will not be scared away by them, but if they come near to bark or snap at him, with his tail, which is long like a whip, he will so slash them, that they will run away and howl."
Reproduction
The Asian water monitor breeds between April and October. The females will lay their eggs about a month after mating in rotting logs or stumps. A clutch can vary from 10 to 40 eggs with an incubation period of 6 to 7 months. When hatched, hatchlings are fully developed and independent. Once males and females reach a length of about 1 meter and 50 centimeters, respectively, they will become reproductively mature and able to breed.
Diet
They are carnivores, and consume a wide range of prey. They are known to eat fish, frogs, rodents, birds, crabs, and snakes. They have also been known to eat turtles as well as young crocodiles and crocodile eggs. Water monitors have been observed eating catfish in a fashion similar to a mammalian carnivore, tearing off chunks of meat with their sharp teeth while holding it with their front legs and then separating different parts of the fish for sequential consumption. In Java, they have also been recorded entering caves at night to hunt bats that have fallen from cave's ceiling.
As carnivores and scavengers, the diet of the Asian water monitor in an urban area in central Thailand includes fish, crabs, Malayan snail-eating turtles (Malayemys macrocephala), Chinese edible frogs (Hoplobatrachus rugulosus), birds, small rodents, domestic cats (Felis catus) and dogs (Canis familiaris), chickens (Gallus gallus domesticus), food scraps and carcass. They are known to feed on dead human bodies. While on the one hand their presence can be helpful in locating a missing person in forensic investigations, on the other hand they can inflict further injuries to the corpse, complicating ascertainment of the cause of death. The stomachs of 20 adult Asian water monitors caught on Redang Island contained mostly human food waste, followed by turtle eggs and hatchlings, crabs and lizard eggs. The monitor does not thrive in these areas, but manage to still live in them. Studies are being conducted in order to understand how these creatures are able to do so in and around human civilization.
In Sri Lanka, human corpses are often scavenged on by V. s. salvator, which can make it hard to identify the deceased, or to run autopsies. For instance, the feeding marks made by a monitor's sharp claws resemble wounds made by bladed weapons. In one case however, the presence of eight dead water monitors near the corpse of a partially scavenged 51 year old man prompted investigation that revealed the possibility that the man died from poisoning after ingesting a bottle of Carbosulfan pesticide, which then poisoned the water monitors that scavenged on his body.
Venom
The possibility of venom in the genus Varanus is widely debated. Previously, venom was thought to be unique to Serpentes (snakes) and Heloderma (venomous lizards). The aftereffects of a Varanus bite were thought to be due to oral bacteria alone, but recent studies have shown venom glands are likely to be present in the mouths of several, if not all, of the species. The venom may be used as a defensive mechanism to fend off predators, to help digest food, to sustain oral hygiene, and possibly to help in capturing and killing prey.
Predation
Adult water monitors have very few predators; with the exception of human hunters, only saltwater crocodiles (Crocodylus porosus) are known to target them.
Interaction with humans
When feeling threatened, water monitors have been known to attack humans, and they should be handled with caution, as their can inflict a severe injury. Still, water monitors have been successfully tamed as pets, and their bites are not known to be fatal.
In 1999, a seven-year-old boy in Pahang, Malaysia, was bitten in the leg while bathing, requiring 18 stitches.
Threats
Monitor lizards are traded globally and are the most common type of lizard to be exported from Southeast Asia, with 8.1 million exported between 1998 and 2007 for the international leather market. Today the majority of the harvesting of feral water monitors occurs in Southeast Asia, in Indonesia, and in peninsular Malaysia. Efforts to breed or farm Water monitors in captivity on a commercial scale have not been widely successful. The Asian water monitor is one of the most exploited varanids; its skin is used for fashion accessories such as shoes, belts and handbags which are shipped globally, with as many as 1.5 million skins traded annually and between 50,000 and 120,000 skins harvested from the wild in peninsular Malaysia. Other uses include a perceived remedy for skin ailments and eczema, novelty food in Indonesia, and a perceived aphrodisiac, and as pets. In India, several tribal communities hunt these monitor lizards for their meat, fat and skin and the eggs are also harvested. They are often considered as pests and their populations are also threatened by habitat loss and habitat fragmentation.
Conservation
In Nepal, it is a protected species under the Wild Animals Protection Act of 2002. In Hong Kong, it is a protected species under Wild Animals Protection Ordinance Cap 170. In Malaysia, this species is one of the most common wild animals, with numbers comparable to the population of macaques there. Although many fall victim to humans via roadkill and animal cruelty, they still thrive in most states of Malaysia, especially in the shrubs of the east coast states such as Pahang and Terengganu as the regulations in Malaysian states differ based on wildlife management authorities. In Thailand, all monitor lizards are protected species. It is still common in large urban areas in Thailand and is frequently seen in Bangkok's canals and parks. Because of this, it is currently listed as Least Concern in the IUCN Red List. These classifications have been made on the basis that this species maintains a geographically wide distribution, can be found in a variety of habitats, adapts to habitats disturbed by humans, and is abundant in portions of its range despite large levels of harvesting.
Loss of habitat and hunting has exterminated water monitors from most of mainland India. In other areas they survive despite being hunted, due in part to the fact that larger ones, including large females that breed large numbers of eggs, have tough skins that are not desirable.
In Sri Lanka, it is protected by local people who value its predation of "crabs that would otherwise undermine the banks of rice fields". It is also protected as it eats venomous snakes.
The species is listed in Appendix II of the Convention on International Trade in Endangered Species (CITES) meaning international trade (import/export) in specimens (including parts and derivatives) is regulated.
| Biology and health sciences | Lizards and other Squamata | Animals |
3741564 | https://en.wikipedia.org/wiki/Type%20%28model%20theory%29 | Type (model theory) | In model theory and related areas of mathematics, a type is an object that describes how a (real or possible) element or finite collection of elements in a mathematical structure might behave. More precisely, it is a set of first-order formulas in a language L with free variables x1, x2,..., xn that are true of a set of n-tuples of an L-structure . Depending on the context, types can be complete or partial and they may use a fixed set of constants, A, from the structure . The question of which types represent actual elements of leads to the ideas of saturated models and omitting types.
Formal definition
Consider a structure for a language L. Let M be the universe of the structure. For every A ⊆ M, let L(A) be the language obtained from L by adding a constant ca for every a ∈ A. In other words,
A 1-type (of ) over A is a set p(x) of formulas in L(A) with at most one free variable x (therefore 1-type) such that for every finite subset p0(x) ⊆ p(x) there is some b ∈ M, depending on p0(x), with (i.e. all formulas in p0(x) are true in when x is replaced by b).
Similarly an n-type (of ) over A is defined to be a set p(x1,...,xn) = p(x) of formulas in L(A), each having its free variables occurring only among the given n free variables x1,...,xn, such that for every finite subset p0(x) ⊆ p(x) there are some elements b1,...,bn ∈ M with .
A complete type of over A is one that is maximal with respect to inclusion. Equivalently, for every either or . Any non-complete type is called a partial type.
So, the word type in general refers to any n-type, partial or complete, over any chosen set of parameters (possibly the empty set).
An n-type p(x) is said to be realized in if there is an element b ∈ Mn such that . The existence of such a realization is guaranteed for any type by the compactness theorem, although the realization might take place in some elementary extension of , rather than in itself.
If a complete type is realized by b in , then the type is typically denoted and referred to as the complete type of b over A.
A type p(x) is said to be isolated by , for , if for all we have . Since finite subsets of a type are always realized in , there is always an element b ∈ Mn such that φ(b) is true in ; i.e. , thus b realizes the entire isolated type. So isolated types will be realized in every elementary substructure or extension. Because of this, isolated types can never be omitted (see below).
A model that realizes the maximum possible variety of types is called a saturated model, and the ultrapower construction provides one way of producing saturated models.
Examples of types
Consider the language L with one binary relation symbol, which we denote as . Let be the structure for this language, which is the ordinal with its standard well-ordering. Let denote the first-order theory of .
Consider the set of L(ω)-formulas . First, we claim this is a type. Let be a finite subset of . We need to find a that satisfies all the formulas in . Well, we can just take the successor of the largest ordinal mentioned in the set of formulas . Then this will clearly contain all the ordinals mentioned in . Thus we have that is a type.
Next, note that is not realized in . For, if it were there would be some that contains every element of .
If we wanted to realize the type, we might be tempted to consider the structure , which is indeed an extension of that realizes the type. Unfortunately, this extension is not elementary, for example, it does not satisfy . In particular, the sentence is satisfied by this structure and not by .
So, we wish to realize the type in an elementary extension. We can do this by defining a new L-structure, which we will denote . The domain of the structure will be where is the set of integers adorned in such a way that . Let denote the usual order of . We interpret the symbol in our new structure by . The idea being that we are adding a "-chain", or copy of the integers, above all the finite ordinals. Clearly any element of realizes the type . Moreover, one can verify that this extension is elementary.
Another example: the complete type of the number 2 over the empty set, considered as a member of the natural numbers, would be the set of all first-order statements (in the language of Peano arithmetic), describing a variable x, that are true when x = 2. This set would include formulas such as , , and . This is an example of an isolated type, since, working over the theory of the naturals, the formula implies all other formulas that are true about the number 2.
As a further example, the statements
and
describing the square root of 2 are consistent with the axioms of ordered fields, and can be extended to a complete type. This type is not realized in the ordered field of rational numbers, but is realized in the ordered field of reals. Similarly, the infinite set of formulas (over the empty set) {x>1, x>1+1, x>1+1+1, ...} is not realized in the ordered field of real numbers, but is realized in the ordered field of hyperreals. Similarly, we can specify a type that is realized by an infinitesimal hyperreal that violates the Archimedean property.
The reason it is useful to restrict the parameters to a certain subset of the model is that it helps to distinguish the types that can be satisfied from those that cannot. For example, using the entire set of real numbers as parameters one could generate an uncountably infinite set of formulas like , , ... that would explicitly rule out every possible real value for x, and therefore could never be realized within the real numbers.
Stone spaces
It is useful to consider the set of complete n-types over A as a topological space. Consider the following equivalence relation on formulas in the free variables x1,..., xn with parameters in A:
One can show that if and only if they are contained in exactly the same complete types.
The set of formulas in free variables x1,...,xn over A up to this equivalence relation is a Boolean algebra (and is canonically isomorphic to the set of A-definable subsets of Mn). The complete n-types correspond to ultrafilters of this Boolean algebra. The set of complete n-types can be made into a topological space by taking the sets of types containing a given formula as a basis of open sets. This constructs the Stone space associated to the Boolean algebra, which is a compact, Hausdorff, and totally disconnected space.
Example. The complete theory of algebraically closed fields of characteristic 0 has quantifier elimination, which allows one to show that the possible complete 1-types (over the empty set) correspond to:
Roots of a given irreducible non-constant polynomial over the rationals with leading coefficient 1. For example, the type of square roots of 2. Each of these types is an isolated point of the Stone space.
Transcendental elements, which are not roots of any non-zero polynomial. This type is a point in the Stone space that is closed but not isolated.
In other words, the 1-types correspond exactly to the prime ideals of the polynomial ring Q[x] over the rationals Q: if r is an element of the model of type p, then the ideal corresponding to p is the set of polynomials with r as a root (which is only the zero polynomial if r is transcendental). More generally, the complete n-types correspond to the prime ideals of the polynomial ring Q[x1,...,xn], in other words to the points of the prime spectrum of this ring. (The Stone space topology can in fact be viewed as the Zariski topology of a Boolean ring induced in a natural way from the Boolean algebra. While the Zariski topology is not in general Hausdorff, it is in the case of Boolean rings.) For example, if q(x,y) is an irreducible polynomial in two variables, there is a 2-type whose realizations are (informally) pairs (x,y) of elements with q(x,y)=0.
Omitting types theorem
Given a complete n-type p one can ask if there is a model of the theory that omits p, in other words there is no n-tuple in the model that realizes p.
If p is an isolated point in the Stone space, i.e. if {p} is an open set, it is easy to see that every model realizes p (at least if the theory is complete). The omitting types theorem says that conversely if p is not isolated then there is a countable model omitting p (provided that the language is countable).
Example: In the theory of algebraically closed fields of characteristic 0, there is a 1-type represented by elements that are transcendental over the prime field Q. This is a non-isolated point of the Stone space (in fact, the only non-isolated point). The field of algebraic numbers is a model omitting this type, and the algebraic closure of any
transcendental extension of the rationals is a model realizing this type.
All the other types are "algebraic numbers" (more precisely, they are the sets of first-order statements satisfied by some given algebraic number), and all such types are realized in all algebraically closed fields of characteristic 0.
| Mathematics | Model theory | null |
6670017 | https://en.wikipedia.org/wiki/Bornean%20orangutan | Bornean orangutan | The Bornean orangutan (Pongo pygmaeus) is an orangutan species endemic to the island of Borneo. It belongs to the only genus of great apes native to Asia and is the largest of the three Pongo species. It has a coarse, reddish coat and up to long arms. It is sexually dimorphic, with females ranging in body length from and males from .
The Bornean orangutan inhabits Borneo lowland rain forests and Borneo montane rain forests up to an elevation of . Its diet includes fruits, seeds, flowers, bird eggs, sap and vines. It is highly intelligent, displaying tool use and distinct cultural patterns. It is critically endangered, with deforestation, palm oil plantations, and hunting posing a serious threat to its survival.
Taxonomy
The Bornean orangutan and the Sumatran orangutan diverged about 400,000 years ago, with a continued low level of gene flow between them since then. The two orangutan species were considered merely subspecies until 1996; they were elevated to species following sequencing of their mitochondrial DNA.
The Bornean orangutan has three subspecies:
Northwest Bornean orangutan Pongo pygmaeus pygmaeus – Sarawak (Malaysia) and northern West Kalimantan (Indonesia).
Central Bornean orangutan P. p. wurmbii – Southern West Kalimantan and Central Kalimantan (Indonesia).
Northeast Bornean orangutan P. p. morio – East Kalimantan (Indonesia) and Sabah (Malaysia).
There is some uncertainty about this, however. The population currently listed as P. p. wurmbii may be closer to the Sumatran orangutan (P. abelii) than to the Bornean orangutan. If this is confirmed, P. abelii would be a subspecies of P. wurmbii (Tiedeman, 1808). In addition, the type locality of P. pygmaeus has not been established beyond doubt; it may be from the population currently listed as P. wurmbii (in which case P. wurmbii would be a junior synonym of P. pygmaeus, while one of the names currently considered a junior synonym of P. pygmaeus would take precedence for the taxon in Sarawak and northern West Kalimantan). Bradon-Jones et al. considered P. morio to be a synonym of P. pygmaeus, and the population found in East Kalimantan and Sabah to be a potentially unnamed separate taxon.
In early October 2014, researchers from domestic and foreign countries found about 50 orangutans in several groups in South Kalimantan Province, although previously there is no record that the province has orangutans.
As a member of the family Hominidae, Bornean orangutans are one of the closest extant relatives to Homo sapiens.
This species was originally discovered by native Malaysians. There are several mentions of orangutans in their folklore. However, this species was originally named and described by the notable zoologist Carl Linnaeus in 1799. Its original name was Simia satyrus, meaning "satyr monkey", but was changed when scientists discovered that not all orangutans are one species. The holotype of this organism is located in the British Museum in London.
The current species name P. pygmaeus is not Latin unlike most other Linnean classifications. The genus name Pongo is derived from the Bantu word used to indicate a large primate. It was originally used to describe chimpanzees in Western African dialects. The species name pygmaeus is derived from the Greek word "pygmy" meaning dwarf.
Description
The Bornean orangutan is the third-largest ape after the western gorilla, and the largest truly arboreal (or tree-dwelling) extant ape. Body weights broadly overlap with the considerably taller Homo sapiens, but the latter is considerably more variable in size. By comparison, the Sumatran orangutan is similar in size but, on average, is marginally lighter in weight. A survey of wild orangutans found that males weigh on average , ranging from , and long; females average , ranging from , and long. While in captivity, orangutans can grow considerably overweight, up to more than . The heaviest known male orangutan in captivity was an obese male named "Andy", who weighed in 1959 when he was 13 years old.
The Bornean orangutan has a distinctive body shape with very long arms that may reach up to 1.5 metres in length. It has grey skin, a coarse, shaggy, reddish coat and prehensile, grasping hands and feet. Its coat does not cover its face unlike most mammals, although Bornean orangutans do have some hair on their faces including a beard and mustache. It also has large, fatty cheek pads known as flanges as well as a pendulous throat sac.
Bornean orangutans are highly sexually dimorphic and have several features that differ between males and females. Males have much larger cheek pads, or flanges, that are composed of muscle and large amounts of fat. In females, the flanges are mostly composed of muscle. Males have relatively larger canines and premolars. Males have a more pronounced beard and mustache. The throat sac in males is also considerably larger. There are two body types for sexually mature males: smaller or larger. Larger males are more dominant but smaller males still breed successfully. There is little sexual dimorphism at birth.
Distribution and habitat
The Bornean orangutan lives in tropical rain forests in the Bornean lowlands, as well as montane rain forests in mountainous areas up to . This species lives throughout the canopy of primary and secondary forest, and moves large distances to find trees bearing fruit.
It is found in the two Malaysian states of Sabah and Sarawak, and four of the five Indonesian Provinces of Kalimantan. Due to habitat destruction, the species distribution is now highly patchy throughout the island, the species has become rare in the southeast of the island, as well as in the forest between the Rajang River in central Sarawak and the Padas River in western Sabah. Its presence in Brunei is uncertain and unconfirmed.
The first complete orangutan skeleton that was discovered was in the Hoa Binh province in Vietnam and thought to be from the late Pleistocene epoch. It differed from modern orangutans only in that its body was proportionately smaller compared to its head. This fossil and others confirm that orangutans once inhabited continental Southeast Asia even though currently, Bornean orangutans are only found in Malaysia and Indonesia.
Behavior and ecology
In history, orangutans ranged throughout Southeast Asia and into southern China, as well as on the island of Java and in southern Sumatra. They primarily inhabit peat swamp forest, tropical heath forest, and mixed dipterocarp forest.
Bornean orangutan are more solitary than their Sumatran relatives. Two or three orangutans with overlapping territories may interact, but only for short periods of time. Although orangutans are not territorial, adult males will display threatening behaviors upon meeting other males, and only socialize with females to mate. Males are considered the most solitary of the orangutans. The Bornean orangutan has a lifespan of 35–45 years in the wild; in captivity it can live to be about 60.
Despite being arboreal, the Bornean orangutan travels on the ground more than its Sumatran counterpart. This may be in part because no large terrestrial predators could threaten an orangutan in Borneo. In Sumatra, orangutans must face predation by the fierce Sumatran tiger.
The Bornean orangutan exhibits nest-building behavior. Nests are built for use at night or during the day. Young orangutans learn by observing their mother's nest-building behaviour. This skill is practiced by juvenile orangutans. Nests may be elaborate and involve a foundation and mattress made by intertwining leaves and branches and adding broken leafy branches. Additional features such as shade, waterproof roof, "pillow", and "blanket", all of which are made from branches, twigs and leaves, may also be added. Nest-building in primates is considered as an example of tool use and not animal architecture.
Diet
The Bornean orangutan diet is composed of over 400 types of food, including wild figs, durians (Durio zibethinus and D. graveolens), leaves, seeds, bird eggs, flowers, sap, vines, honey, fungi, spider webs, insects, and, to a lesser extent than the Sumatran orangutan, bark. They have also been known to consume the inner shoots of plants and vines. They will also occasionally eat nutrient rich soil. They get the necessary quantities of water from both fruit and from tree holes.
Bornean orangutans have been sighted using spears to attempt (unsuccessfully) to catch fish. The species has been observed using tools such as leaves to wipe off faeces, a pad of leaves for holding spiny durian fruit, a leafy branch for a bee swatter, a bunch of leafy branches held together as an "umbrella" while traveling in the rain, a single stick as backscratcher, and a branch or tree trunk as a missile. In some regions, orangutans occasionally eat soil to get minerals that may neutralize the toxins and acids they consume in their primarily vegetarian diets. On rare occasions, orangutans will prey upon other, smaller primates, such as slow lorises.
Reproduction
Males and females generally come together only to mate. Subadult males try to mate with any female and are successful around half the time. Dominant flanged males will call and advertise their position to receptive females, who prefer mating with flanged males. Adult males will often target females with weaned infants as mating partners because the female is likely to be fertile.
Females reach sexual maturity and experience their first ovulatory cycle between about six and 11 years of age, although females with more body fat may experience this at an earlier age. The estrous cycle lasts between 22 and 30 days and menopause has been reported in captive orangutans at about age 48. Females tend to give birth at about 14–15 years of age. Newborn orangutans nurse every three to four hours, and begin to take soft food from their mothers' lips by four months. During the first year of its life, the young clings to its mother's abdomen by entwining its fingers in and gripping her hair. Offspring are weaned at about four years, but this could be much longer, and soon after they start their adolescent stage of exploring, but always within sight of their mother. During this period, they will also actively seek other young orangutans to play with and travel with. On average, juveniles do not become completely independent until they are about seven years of age. The birth rate for orangutans has been decreasing largely due to a lack of sufficient nutrients as a result of habitat loss.
A 2011 study on female orangutans in free-ranging rehabilitation programs found that individuals that were supplemented with food resources had shorter interbirth intervals, as well as a reduced age, at first birth.
Conservation status
The Bornean orangutan is more common than the Sumatran, with about 104,700 individuals in the wild, whereas just under 14,000 Sumatran orangutans are left in the wild. Orangutans are becoming increasingly endangered due to habitat destruction and the bushmeat trade, and young orangutans are captured to be sold as pets, usually entailing the killing of their mothers.
The Bornean orangutan is critically endangered according to the IUCN Red List of mammals, and is listed on Appendix I of CITES. The total number of Bornean orangutans is estimated to be less than 14% of what it was in the recent past (from around 10,000 years ago until the middle of the 20th century), and this sharp decline has occurred mostly over the past few decades due to human activities and development. Species distribution is now highly patchy throughout Borneo; it is apparently absent or uncommon in the southeast of the island, as well as in the forest between the Rajang River in central Sarawak and the Padas River in western Sabah (including the Sultanate of Brunei). A population of around 6,900 is found in Sabangau National Park, but this environment is at risk.
This view is also supported by the United Nations Environment Programme, which stated in its 2007 report that, due to illegal logging, fire and the extensive development of palm oil plantations, orangutans are critically endangered, and if the current trend continues, they will become extinct. When forest is burned down to clear room for palm oil plantations, not only does the Bornean orangutan suffer from habitat loss, but several individuals have been burned and killed in fires. Palm oil accounts for over one tenth of Indonesia's export earnings. It is in high demand because it is used in several packaged foods, deodorants, shampoos, soaps, candies, and baked goods.
Climate change is another threat to Bornean orangutan conservation. The effects that human activity have had on Indonesian rainfall have made food less abundant and so Bornean orangutans are less likely to receive full nutrients so that they can be sufficiently healthy to breed.
A November 2011 survey, based on interviews with 6,983 respondents in 687 villages across Kalimantan in 2008 to 2009, gave estimated orangutan killing rates of between 750 and 1800 in the year leading up to April 2008. These killing rates were higher than previously thought and confirm that the continued existence of the orangutan in Kalimantan is under serious threat. The survey did not quantify the additional threat to the species due to habitat loss from deforestation and expanding palm-oil plantations. The survey found that 73% of respondents knew orangutans were protected by Indonesian law.
However, the Indonesian government rarely prosecutes or punishes perpetrators. In a rare prosecution in November 2011, two men were arrested for killing at least 20 orangutans and a number of long-nosed proboscis monkeys. They were ordered to conduct the killings by the supervisor of a palm oil plantation, to protect the crop, with a payment of $100 for a dead orangutan and $22 for a monkey.
Rescue and rehabilitation centers
A number of orangutan rescue and rehabilitation projects operate in Borneo.
The Borneo Orangutan Survival Foundation (BOS) founded by Willie Smits has rescue and rehabilitation centres at Wanariset and Samboja Lestari in East Kalimantan and Nyaru Menteng, in Central Kalimantan founded and managed by Lone Drøscher Nielsen. BOS also works to conserve and recreate the fast-disappearing rainforest habitat of the orangutan, at Samboja Lestari and Mawas.
Orangutan Foundation International, founded by Birutė Galdikas, rescues and rehabilitates orangutans, preparing them for release back into protected areas of the Indonesian rain forest. In addition, it promotes the preservation of the rain forest for them.
The Sepilok Orangutan Rehabilitation Centre near Sandakan in the state of Sabah in Malaysian Borneo opened in 1964 as the first official orangutan rehabilitation project.
Orangutan Foundation, founded by Ashley Leiman, operates programmes in Central Kalimantan, Indonesian Borneo. The Foundation rescues orphaned orangutans and enters them into their soft-release programme, allowing them to develop the skills necessary to survive in the wild. When old enough, orangutans are released into the protected Lamandau Wildlife Reserve. Orangutan Foundation works to protect orangutans by focusing on habitat protection and capacity building, especially in local communities.
A seven-year longitudinal study published in 2011 looked at whether the lifespan of zoo-housed orangutans was related to a subjective assessment of well-being, with the intent of applying such measures to assess the welfare of orangutans in captivity. Of the subjects, 100 were Sumatran (Pongo abelii), 54 Bornean (Pongo pygmaeus) and 30 were hybrid orangutans. 113 zoo employees, who were highly familiar with the typical behavior of the orangutans, used a four-item questionnaire to assess their subjective well-being. The results indicated that orangutans in higher subjective well-being were less likely to die during the follow-up period. The study concluded that happiness was related to longer life in orangutans.
In late 2014, Nyaru Menteng veterinarians failed to rescue the life of a female orangutan. An operation was performed in which 40 air-rifle pellets were removed from her body. The orangutan was found at a palm oil plantation in Indonesian Borneo.
Genome and demographic history
Orangutans and humans diverged lineages approximately 14–18 million years ago. About 17,000 years ago, there was a migration of the Bornean orangutans as they eventually went to Sumatra, effectively trading places with the Sumatra orangutans that were there at the time. These two species of orangutans have been closely related throughout their evolutionary history due to the fact that they were so close in physical proximity. Therefore, their genomes and demographic history are similar. The two species themselves are estimated to have split about 3.5 million years ago. Although these two species have officially diverged, it is speculated that the reason as to why they are genetically similar is because the males of each respective species tend to migrate between the two islands and breed with the females from their sister species. As a result, both the Bornean orangutans and the Sumatran orangutans have been studied closely as a pair, and thus much genome findings attribute evolutionary changes to this relationship. In addition, the Bornean orangutans, as compared to the Sumatran orangutans, have lower autosomal gene diversity. This is attributed to the fact that they have a much smaller population size. Also, the Bornean orangutans have lower nucleotide diversity.
As the Bornean orangutans and Sumatra orangutans both exist within the same genus, they exhibit similar cultural behaviors that have been found to exist amongst most orangutan populations. The fact that orangutans tend to showcase similar cultural traditions is due to the fact that they typically live in similar environments and are adept at learning from one another from their early stages of life.
The Bornean orangutan has been linked to the fact that it has gone through a deep divergence in relation to its relatives and ancestors. During the Middle Pleistocene, there were low levels of gene flow, which was determined through the analysis of Y-chromosomal data. One reason as to why this may have occurred is because of the Sunda shelf, which is where the island of Borneo is located. During this time, this event's dry climate during the Late Pleistocene attributed to a more abundant genetic exchange. As a result, there were many early divergences of gene pools between the Bornean orangutans, as well as the Sumatran orangutans. Relating back to the Middle Pleistocene, the Bornean orangutan lineage went through a dramatic population decline. This is likely attributed to the fact that they had been isolated from their ancestral populations. Therefore, natural geographic barriers are attributed to be the reason as to why the Bornean orangutans were eventually isolated and ended up colonizing other regions. In addition, this geographic isolation also indicates that the Bornean orangutans did not undergo a severe genetic bottleneck. With the Borneo orangutan, selection was found to have been found through physiological adaptations – most of which has to do with being able to adapt to the ever-changing climate on the Borneo island.
| Biology and health sciences | Apes | Animals |
1359660 | https://en.wikipedia.org/wiki/Giant%20pouched%20rat | Giant pouched rat | The giant pouched rats (genus Cricetomys) of sub-Saharan Africa are large muroid rodents.
Description
Their head and body lengths range from with scaly tails ranging from . They weigh between .
Taxonomy
Giant pouched rats are only distantly related to the true rats, although until recently they had been placed in the same family, Muridae. Recent molecular studies, however, place them in the family Nesomyidae, part of an ancient radiation of African and Malagasy muroids. The name "pouched rat" refers to their large cheek pouches.
The species are:
Behaviour
Females have been said to be capable of producing up to 10 litters yearly. Gestation is 27–36 days. The animals generally have between six and eight nipples. One to five young are born at a time.
The animals are nocturnal omnivores, and feed on vegetation and small animals, especially insects. They have a particular taste for palm nuts.
Interaction with humans
In many African countries, giant pouched rats are valued as an important food item.
They are easily tamed as pets, but were associated with an outbreak of monkeypox in the USA in 2003, and have since been banned from importation to the U.S.
Detecting explosives and tuberculosis by scent
These rats are also becoming useful in some areas for detecting land mines; their acute sense of smell is very effective in detecting explosives such as TNT, and at the same time they are light enough to not trigger any of the mines including antipersonnel mines. The rats are being trained by APOPO, a nonprofit social venture based in Tanzania.
The procedure for training rats to detect land mines was conceived of and developed by Belgian Bart Weetjens. Training starts at four weeks of age, when the rats are handled to accustom them to humans and exposed to a variety of sights and sounds. They learn to associate a clicker with a food reward of banana or banana-peanut paste. They are then trained to indicate a hole that contains TNT by nosing it for five seconds. Then they learn to find the correct hole in a line of holes. Finally, the rat is trained to wear a harness and practises outdoors on a lead, finding inactive mines under soil. At the end of their training, they are tested; they must find all the mines in an area of that has been seeded with inactivated mines. It is a blind test; their handlers do not know where the mines are. If they succeed, they are certified as bomb-sniffing rats.
APOPO is also training the rats to detect tuberculosis by sniffing sputum samples to detect Mycobacterium tuberculosis. The rats can test many more samples than laboratory techniques can—100 in 20 minutes, which would take a lab technician up to four days using conventional microscopy. Furthermore, samples submitted for secondary screening by the rats reassess 52% of initially negative tests are as positive. In some cases TB detected by rats has not been confirmed by clinical tests, but patients later developed TB, suggesting that rats can detect the disease before a clinical test. they were being used to screen for tuberculosis in Tanzania and Ethiopia.
Land mine and tuberculosis sniffing rats are called HeroRATs.
In popular culture
Ben, in the 2003 remake of Willard, was a Gambian pouched rat.
| Biology and health sciences | Rodents | Animals |
1359960 | https://en.wikipedia.org/wiki/Pomology | Pomology | Pomology (from Latin , "fruit", + , "study") is a branch of botany that studies fruits and their cultivation. Someone who researches and practices the science of pomology is called a pomologist. The term fruticulture (from Latin , "fruit", + , "care") is also used to describe the agricultural practice of growing fruits in orchards.
Pomological research is mainly focused on the development, enhancement, cultivation, and physiological studies of fruit trees. The goals of fruit tree improvement include enhancement of fruit quality, regulation of production periods, and reduction of production costs.
History
Middle East
In ancient Mesopotamia, pomology was practiced by the Sumerians, who are known to have grown various types of fruit, including dates, grapes, apples, melons, and figs. While the first fruits cultivated by the Egyptians were likely indigenous, such as the palm date and sorghum, more fruits were introduced as other cultural influences were introduced. Grapes and watermelon were found throughout predynastic Egyptian sites, as were the sycamore fig, dom palm, and Christ's thorn. The carob, olive, apple, and pomegranate were introduced to Egyptians during the New Kingdom. Later, during the Greco-Roman period peaches and pears were also introduced.
Europe
The ancient Greeks and Romans also had a healthy tradition of pomology, and they cultivated a wide range of fruits, including apples, pears, figs, grapes, quinces, citron, strawberries, blackberries, elderberries, currants, damson plums, dates, melons, rose hips, and pomegranates. Less common fruits were the more exotic azeroles and medlars. Cherries and apricots, both introduced in the 1st century BC, were popular. Peaches were introduced in the 1st century AD from Persia. Oranges and lemons were known but used more for medicinal purposes than in cookery. The Romans, in particular, were known for their advanced methods of fruit cultivation and storage, and they developed many of the approaches that are still used in modern pomology.
United States
During the mid-19th century in the United States, farmers were expanding fruit orchard programs in response to growing markets. At the same time, horticulturists from the USDA and agricultural colleges were bringing new varieties to the US from foreign expeditions, and developing experimental lots for these fruits. In response to this increased interest and activity, the USDA established the Division of Pomology in 1886 and named Henry E. Van Deman as chief pomologist. An important focus of the division was to publish illustrated accounts of new varieties and to disseminate research findings to fruit growers and breeders through special publications and annual reports. During this period Andrew Jackson Downing and his brother Charles were prominent in pomology and horticulture, producing The Fruits and Fruit Trees of America (1845).
The introduction of new varieties required an exact portrait of the fruit so that plant breeders could accurately document and disseminate their research results. Since the use of scientific photography was not widespread in the late 19th century, the USDA commissioned artists to create watercolor illustrations of newly introduced cultivars. Many of the watercolors were used for lithographic reproductions in USDA publications, such as the Report of the Pomologist and the Yearbook of Agriculture. Today, the collection of approximately 7,700 watercolors is preserved in the National Agricultural Library's Special Collections, where it serves as a major historic and botanic resource to a variety of researchers, including horticulturists, historians, artists, and publishers.
| Technology | Horticulture | null |
1360138 | https://en.wikipedia.org/wiki/Corythosaurus | Corythosaurus | Corythosaurus (; ) is a genus of hadrosaurid "duck-billed" dinosaur from the Late Cretaceous period, about 77–75.7 million years ago, in what is now western North America. Its name is derived from the Greek word κόρυς, meaning "helmet", named and described in 1914 by Barnum Brown. Corythosaurus is now thought to be a lambeosaurine, thus related to Lambeosaurus, Nipponosaurus, Velafrons, Hypacrosaurus, and Olorotitan. Corythosaurus has an estimated length of and has a skull, including the crest, that is tall.
Corythosaurus is known from many complete specimens, including the nearly complete holotype found by Brown in 1911. The holotype skeleton is only missing the last section of the tail and part of the front legs, but was preserved with impressions of polygonal scales. Corythosaurus is known from many skulls with tall crests that resemble those of the cassowary and a Corinthian helmet. The most likely function of the crest is thought to be vocalization. As in a trombone, sound waves would travel through many chambers in the crest and then get amplified when Corythosaurus exhaled. One Corythosaurus specimen has even been preserved with its last meal in its chest cavity. Inside the cavity were remains of conifer needles, seeds, twigs, and fruits, suggesting that Corythosaurus probably fed on all of these.
The two species of Corythosaurus are both present in slightly different levels of the Dinosaur Park Formation. Both still co-existed with theropods and other ornithischians, like Daspletosaurus, Brachylophosaurus, Parasaurolophus, Scolosaurus, and Chasmosaurus.
Discovery and species
The first specimen, AMNH 5240, was discovered in 1911 by Barnum Brown in Red Deer River of Alberta and secured by him in the Fall of 1912. As well as an almost complete skeleton, the find was notable because impressions of much of the creature's skin had also survived. The specimen came from the Belly River Group of the province. The left or underside of the skeleton was preserved in carbonaceous clay, making it difficult to expose the skin to the elements. The skeleton was articulated and only missing about the last of the tail and front legs. Both scapulae and coracoids are preserved in position, but the rest of the front legs are gone (except for phalanges and pieces of the humeri, ulnae, and radii). Apparently, the remaining front legs were weathered or eroded away. Impressions of the integument were preserved covering over a large part of the skeleton’s outlining and shows the form of the body. Another specimen, AMNH 5338, was found in 1914 by Brown and Peter Kaisen. Both specimens are now housed in the American Museum of Natural History in their original death poses.
The type species, Corythosaurus casuarius, was named by Barnum Brown in 1914, based on the first specimen collected by him in 1912. AMNH 5240 is thus the holotype. In 1916, the original author, Brown, published a more detailed description that was also based on AMNH 5338, which is therefore the plesiotype. Corythosaurus is among many lambeosaurines that possess crests and it was the crest that lends Corythosaurus its name. The generic name Corythosaurus is derived from the Greek κόρυθος,(korythos), "Corinthian helmet", and means "helmeted lizard". The specific name, casuarius, refers to the cassowary, a bird with a similar skull crest. The full binomial of Corythosaurus casuarius thus means "Cassowary-like reptile, with a Corinthian helmet crest".
The two best preserved specimens of Corythosaurus, found by Charles H. Sternberg in 1912, were lost on December 6, 1916, while being carried by the SS Mount Temple to the United Kingdom during World War I. They were being sent to Arthur Smith Woodward, a paleontologist of the British Museum of Natural History in England, when the ship transporting them was sunk by the German merchant raider in the middle of the ocean.
There were formerly up to seven species described, including C. casuarius, C. bicristatus (Parks 1935), C. brevicristatus (Parks 1935), C. excavatus (Gilmore 1923), C. frontalis (Parks 1935), and C. intermedius (Parks 1923). In 1975, Peter Dodson studied the differences between the skulls and crests of different species of lambeosaurine dinosaurs. He found that the differences in size and shape may have actually been related to the sex and age of the animal. Only one species is currently recognized for certain, C. casuarius, although C. intermedius has been recognized as valid in some studies. It is based on specimen ROM 776, a skull found by Levi Sternberg in 1920, was named by William Parks in 1923. Originally, he named it Stephanosaurus intermedius earlier that year. The specific name of C. intermedius is derived from its apparent intermediate position according to Parks. C. intermedius lived at a slightly later time in the Campanian than C. casuarius and the two species are not identical, which supported the separation of them in a 2009 study. The invalid species, C. excavatus (specimen UALVP 13), was based on only a skull found in 1920 and wouldn't be reunited with the rest of its remains until 2012.
Description
Size
Benson et al. (2012) estimated that Corythosaurus has an average length of . In 1962, Edwin H. Colbert used models of specific dinosaurs, including Corythosaurus, to estimate their weight. The Corythosaurus model used was modelled by Vincent Fusco, after a mounted skeleton, and supervised by Barnum Brown. After testing, it was concluded that the average weight of Corythosaurus was . The total length of Corythosaurus specimen AMNH 5240 was found to be long, with a weight close to . In 2016, Gregory S. Paul estimated that C. casuarius reached long and in weight and that C. intermedius reached in length and in weight. A "morphologically adult-sized specimen" of C. casuarius measured approximately long.
Proportionally, the skull is much shorter and smaller than that of Edmontosaurus (formerly Trachodon), Kritosaurus, or Saurolophus. But, when including its crest, its superficial area is almost as large.
Skull
Over twenty skulls have been found from this dinosaur. As with other lambeosaurines, the animal bore a tall, elaborate, bony crest atop its skull that contained the elongate narial passages. The narial passages extended into the crest, first into separate pockets in the sides, then into a single central chamber, and onward into the respiratory system. The skull of the type specimen has no dermal impressions on it. During preservation, it was compressed laterally and the width is now about two-thirds of what it would have been in real life. According to Brown, the compression also caused the nasals to shift where they pressed down on the premaxillaries. Because they were pressed on the premaxillaries, the nasals would have closed the nares. Apart from the compression, the skull appears to be normal. Contrary to what Brown assumed, the areas concerned were fully part of the praemaxillae.
As aforementioned, the crests of Corythosaurus resemble that of a cassowary or a Corinthian helmet. They are formed by a combination of the praemaxillae, nasals, prefrontals, and frontals, as in Saurolophus, but instead of projecting backwards as a spine, they rise up to make the highest point above the orbit. The two halves of the crest are separated by a median suture. In front of the orbit, the crest is made of thick bone.
The nasals make up most of the crest. Brown assumed that they extended from the beaks' tip to the highest spot along the crest and that, unlike those in other genera, the nasals meet in the center and are not separated in front by an ascending premaxillary process. However, Brown mistook the praemaxillae for the nasals. The snout is actually largely formed by them and they do separate the nasals. Brown also thought that, on the top and back of the crest, the whole external face is covered by the frontals. Again he made a mistake, as what he assumed to be the frontals are in fact the nasals. The nasals end at the back of the squamosals in a hooked, short process. The prefrontals also make up part of the crest. However, Brown mistook the lower upper branch of the praemaxilla for the prefrontal. The actual prefrontal, which is triangular in shape, is located at the side of the crest base. It was seen by Brown as a part of the frontal. The real frontals, which are largely internal to the crest base structure, are not visible from the side.
The mouth of the holotype of Corythosaurus is narrow. The praemaxillae each form two long folds that enclose air passages extending the narial passages to the front of the snout. There, they end in narrow openings, sometimes called "pseudonares", which are false bony nostrils. These were mistaken by Brown for the real nares or nostrils. These are actually situated inside the crest, above the eye sockets. As in Saurolophus, the expanded portion of the premaxillary in front of the pseudonaris' opening is elongate. By comparison, the bill of Kritosaurus is short and the pseudonares extend far forward. At the end of the Corythosaurus bill, the two pseudonares unite into one. Because of his incorrect identification, Brown assumed that the holotype's inferior process of the premaxillary was shorter than in Kritosaurus and Saurolophus and that the process does not unite with the lacrimal, which is another difference from those genera. The praemaxilla actually does touch the lacrimal and extends to the rear until well behind the eye socket.
The lower jaw of the holotype is long and deep. The total length of the crest from the beak to the uppermost tip of the type specimen is , its total length is , and its height is .
Soft tissue
In the holotype of C. casuarius, the sides and tail are covered in scales of several types. Polygonal tuberculate scales, covered in small bumps, vary in size across the body. Conical limpet-like scales are only preserved on a fold of skin preserved on the back of the tibia, but this was probably from the bottom of the belly instead of the leg. Separating the polygonal scales of C. casuarius are shieldlike scales, arranged close together in rows. Ossified tendons are present on all the vertebrae, except for those in the cervical region. On no vertebrae do the tendons extend below the transverse processes. Each tendon is flattened at its origin, transversely ovoid in the central rod, and ends at a rounded point.
Aside from those found on Corythosaurus casuarius, extensive skin impressions have been found on Edmontosaurus annectens and notable integument has also been found on Brachylophosaurus canadensis, Gryposaurus notabilis, Parasaurolophus walkeri, Lambeosaurus magnicristatus, L. lambei, Saurolophus angustirsotris, and on unidentified ornithopods. Of these, L. lambei, C. casuarius, G. notabilis, P. walkeri, and S. angustirsotris have preserved polygonal scales. The scales on L. lambei, S. angustirostris, and C. casuarius are all similar. Corythosaurus is one of very few hadrosaurids which have preserved skin impressions on the hind limbs and feet. A study in 2013 showed that, amongst hadrosaurids, Saurolophus angustirostris preserved the best and most complete foot and limb integument, although other species like S. osborni, Edmontosaurus annectens, and Lambeosaurus lambei (= L. clavinitialis) share a fair amount of preserved tissue on those regions.
It was once thought that this dinosaur lived mostly in the water, due to the appearance of webbed hands and feet. However, it was later discovered that the so-called "webs" were in fact deflated padding, much like that found on many modern mammals.
Distinguishing characteristics
A set of characters were indicated by Barnum Brown in 1914 to distinguish Corythosaurus from all other hadrosaurids from Alberta. These include a comparatively short skull with a high helmet-like crest formed by the nasals, prefrontals, and frontals; the nasals not being separated in front by the premaxillaries; a narrow beak with an expansion in front of an elongated naris; and a small narial opening.
In 1916, Brown expanded the character set to include even more features. In the revised version, these extra features include a comparatively short skull with a high helmet-like crest formed by nasals, prefrontals, and frontals; the nasals not being separated in front by premaxillaries; a narrow beak; expanded section in front of the elongated nares; a small narial opening; a vertebral formula of 15 cervicals, 19 dorsals, 8 sacrals, and 61+ caudals; possession of dorsal spines of a medium height; high anterior caudal spines; long chevrons; long scapulae that possess a blade of medium width; a radius considerably longer than the humerus; comparatively short metacarpals, an anteriorly decurved ilium; a long ischium with a foot-like terminal expansion; a pubis with an anterior blade that is short and broadly expanded at the end; a femur that is longer than the tibia; the phalanges of pes are short; that the integument over the sides and tail composed of polygonal tuberculate scales without pattern, but graded in size in different parts of the body; and a belly with longitudinal rows of large conical limpet-like scales separated by uniformly large polygonal tubercles. Again, the presumed traits of the snout are incorrect because Brown confused the praemaxillae with the nasal bones and the nasal bones with the frontals. Most of the postcranial traits are today known to be shared with various other lambeosaurines.
Classification
Originally, Brown referred to Corythosaurus as a member of the family Trachodontidae (now Hadrosauridae). Inside Trachodontidae were the subfamilies Trachodontinae and Saurolophinae. Brown classified Hadrosaurus, Trachodon, Claosaurus, and Kritosaurus in Trachodontinae, whereas he classified Corythosaurus, Stephanosaurus, and Saurolophus in Saurolophinae.
Later, Brown revised the phylogeny of Corythosaurus, finding that it was closely related and possibly ancestral to Hypacrosaurus. The only differences he found between them were the development of the vertebrae and the proportions of the legs. During a study of dinosaurian ilia in the 1920s, Alfred Sherwood Romer proposed that the two orders of dinosaurs might have evolved separately and that birds, based on the shape and proportions of their ilia, might truly be specialized ornithischians. He used both Tyrannosaurus and Corythosaurus as a base model to analyze which theory is more likely true. He found that, even though birds are thought of as saurischians, it is very plausible for them to have evolved their specific pelvic musculature and anatomy if they evolved from ornithschians like Corythosaurus. However, even though the pelvic structure of Corythosaurus and other ornithischians does bear a greater superficial resemblance to birds than the saurischian pelvis does, birds are now known to be highly derived maniraptoran theropods.
Corythosaurus is currently classified as a hadrosaurid in the subfamily Lambeosaurinae. It is related to other hadrosaurs such as Hypacrosaurus, Lambeosaurus, and Olorotitan. With the exception of Olorotitan, they all share similar looking skulls and crests. However, research published in 2003 has suggested that even though it possesses a unique crest, Olorotitan is Corythosaurus's closest known relative. Benson et al. (2012) found that Corythosaurus was closely related to Velafrons, Nipponosaurus, and Hypacrosaurus, with them forming a group of fan-crested lambeosaurines.
In 2014, a study including the description of Zhanghenglong was published in the journal PLOS ONE. The study included an almost complete cladogram of hadrosauroid relationships, including Corythosaurus as the most derived lambeosaurine and being the sister taxon to Hypacrosaurus. The below cladogram is a simplified version including only Lambeosaurini.
Paleobiology
Comparisons between the scleral rings of Corythosaurus and modern reptiles suggest that it may have been cathemeral, meaning it was most active throughout the day at short intervals. The sense of hearing in hadrosaurids, specifically such as Lophorhothon, also seems to have been greatly developed because of an elongated lagena. The presence of a thin stapes (an ear bone that is rod-like in reptiles), combined with a large eardrum, implies the existence of a sensitive middle ear. It is possible that hadrosaurid ears are sensitive enough to detect as much sound as a modern crocodilian.
Crest function
The internal structures of the crest of Corythosaurus are quite complex, making possible a call that could be used as a warning or for attracting a mate. Nasal passageways of Corythosaurus, as well as Hypacrosaurus and Lambeosaurus, are S-shaped, with Parasaurolophus only possessing U-shaped tubes. Any vocalization would travel through these elaborate chambers and probably get amplified. Scientists speculate that Corythosaurus could make loud, low pitched cries "like a wind or brass instrument", such as a trombone. The sounds could serve to alert other Corythosaurus to the presence of food or a potential threat from a predator. The nasal passages emit low-frequency sounds when Corythosaurus exhaled. The individual crests would produce different sounds, so it is likely that each species of lambeosaurine would have had a unique sound. However, even though the range for different lambeosaurine nasal passages vary greatly, they all probably made low-pitched sounds. This might be because low sounds (below 400 Hz) travel a set distance in any environment, while higher sounds (above 400 Hz) have a larger spread in the distance travelled.
When they were first described, crested hadrosaurs were thought to be aquatic, an assessment based incorrectly on webbing that is now known to be padding. The theory was that the animals could swim deep in the water and use the crest to store air to breath. However, it has now been proven that the crest did not have any holes in the end and the water pressure at even would be too great for the lungs to be able to inflate.
Growth
Corythosaurus casuarius is one of a few lambeosaurines, along with Lambeosaurus lambei, Hypacrosaurus stebingeri, and H. altispinus, to have had surviving fossilized juveniles assigned to it. Juveniles are harder to assign to species because, at a young age, they lack the distinctive larger crests of adults. As they age, lambeosaurine crests tend to grow and become more prominent come maturity. In the Dinosaur Park Formation, over fifty articulated specimens have been found that come from many different genera. Among them, juveniles are hard to identify at the species level. Earlier, four genera and thirteen species were recognized from the formation's area when paleontologists used differences in size and crest shape to differentiate taxa. The smallest specimens were identified as Tetragonosaurus, now seen as a synonym of Procheneosaurus, and the largest skeletons were called either Corythosaurus or Lambeosaurus. An adult was even identified as Parasaurolophus. Small lambeosaurines from the Horseshoe Canyon Formation were referred to Cheneosaurus.
Corythosaurus started developing its crest when they were half the size of adults, but Parasaurolophus juveniles grew crests when they were only 25% as long as adults. Juvenile Corythosaurus, along with adults, had a premaxilla-nasal fontanelle. Young and adult Corythosaurus are similar to Lambeosaurus and Hypacrosaurus, but dissimilar to Parasaurolophus in that the sutures of the skull are sinuous, not smooth and straight. This feature helps to differentiate Parasaurolophini from Lambeosaurini. Generally, the crests of juveniles of lambeosaurines like Corythosaurus, Lambeosaurus, Hypacrosaurus stebingeri, parasaurolophines like Parasaurolophus, and primitive lambeosaurines like Kazaklambia are quite alike, although other features can be used to distinguish them.
Work by Dodson (1975) recognized that there were many less taxa present in Alberta. Tetragonosaurus was found to be juveniles of Corythosaurus or Lambeosaurus. T. erectofrons was assigned to Corythosaurus based largely on biometric information. The only non-typic specimen of Tetragonosaurus, assigned to T. erectofrons, was later found to be referable to Hypacrosaurus, although the holotype of the species was still found to be assignable to Corythosaurus.
Diet
Corythosaurus was an ornithopod, therefore being a herbivore. Benson et al. (2012) realized that the beak of Corythosaurus was shallow and delicate, concluding that it must have been used to feed upon soft vegetation. Based on the climate of the Late Cretaceous, they guessed that Corythosaurus would have been a selective feeder, eating only the juiciest fruits and youngest leaves. A Corythosaurus specimen has been preserved with its last meal in its chest cavity. Inside the cavity were remains of conifer needles, seeds, twigs, and fruits, meaning that Corythosaurus probably fed on all of these, implying that it was a browser.
Paleoecology
Fossils have been found in the upper Oldman Formation and lower Dinosaur Park Formation of Canada. The Oldman Formation dates to the Campanian, about 77.5 to 76.5 million years ago, and the Dinosaur Park Formation dates from 76.6 to 74.8 million years ago. Corythosaurus lived from ~77–75.7 million years ago. In the Dinosaur Park Formation, C. casuarius lived from 76.6 to 75.9 mya, with C. intermedius living from 75.8 to 75.7 mya. In the Oldman Formation, C. casuarius, the only species of Corythosaurus from the deposits, lived about 77 to 76.5 mya. The holotype specimen was clearly a carcass that had floated up on a beach, as Unio shells, water-worn bones, and a baenid turtle were preserved all around it. Corythosaurus probably lived in a woodland forest and might have occasionally wandered into swampy areas.
A limited fauna is known from the upper section of the Oldman Formation and Corythosaurus casuarius, as well as C. intermedius, are among the taxa. Also from the section of the formation are the theropods Daspletosaurus and Saurornitholestes, the hadrosaurids Brachylophosaurus, Gryposaurus, and Parasaurolophus, the ankylosaurid Scolosaurus, and the ceratopsians Coronosaurus and Chasmosaurus. Other genera are known, but do not persist from the upper section of the formation, therefore not being contemporaries of Corythosaurus.
Corythosaurus casuarius is widespread throughout the lower unit of the Dinosaur Park Formation. In it, Corythosaurus was found to be closely associated with the ceratopsid Centrosaurus apertus. Their associating was found in the Dinosaur Park, Judith River, and Mesaverde formations, as well as the Wind River Basin and the Wheatland County area. Corythosaurus lived alongside numerous other giant herbivores, such as the hadrosaurids Gryposaurus and Parasaurolophus, the ceratopsids Centrosaurus and Chasmosaurus, and the ankylosaurids Scolosaurus, Edmontonia, and Dyoplosaurus in the earliest stages of the formation, Dyoplosaurus, Panoplosaurus, and Euoplocephalus in the middle age, and Euoplocephalus alone in later stages of the formation. Studies of the jaw anatomy and mechanics of these dinosaurs suggests they probably all occupied slightly different ecological niches in order to avoid direct competition for food in such a crowded eco-space. The only large predators known from the same levels of the formation as Corythosaurus are the tyrannosaurids Gorgosaurus libratus and an unnamed species of Daspletosaurus.
Thomas M. Lehman has observed that Corythosaurus hasn't been discovered outside of southern Alberta, even though it is one of the most abundant Judithian dinosaurs in the region. Large herbivores like the hadrosaurs living in North America during the Late Cretaceous had "remarkably small geographic ranges" despite their large body size and high mobility. This restricted distribution strongly contrasts with modern mammalian faunas whose large herbivores' ranges "typical[ly] ... span much of a continent."
| Biology and health sciences | Ornitischians | Animals |
1361581 | https://en.wikipedia.org/wiki/Mechanical%20television | Mechanical television | Mechanical television or mechanical scan television is an obsolete television system that relies on a mechanical scanning device, such as a rotating disk with holes in it or a rotating mirror drum, to scan the scene and generate the video signal, and a similar mechanical device at the receiver to display the picture. This contrasts with vacuum tube electronic television technology, using electron beam scanning methods, for example in cathode-ray tube (CRT) televisions. Subsequently, modern solid-state liquid-crystal displays (LCD) and LED displays are now used to create and display television pictures.
Mechanical-scanning methods were used in the earliest experimental television systems in the 1920s and 1930s. One of the first experimental wireless television transmissions was by Scottish inventor John Logie Baird on October 2, 1925, in London. By 1928 many radio stations were broadcasting experimental television programs using mechanical systems. However the technology never produced images of sufficient quality to become popular with the public. Mechanical-scan systems were largely superseded by electronic-scan technology in the mid-1930s, which was used in the first commercially successful television broadcasts which began in the late 1930s. In the U.S., experimental stations such as W2XAB in New York City began broadcasting mechanical television programs in 1931 but discontinued operations on February 20, 1933, until returning with an all-electronic system in 1939.
A mechanical television receiver was also called a televisor.
History
Early research
The first mechanical raster scanning techniques were developed in the 19th century for facsimile, the transmission of still images by wire. Alexander Bain introduced the facsimile machine in 1843 to 1846. Frederick Bakewell demonstrated a working laboratory version in 1851. The first practical facsimile system, working on telegraph lines, was developed and put into service by Giovanni Caselli from 1856 onward.
Willoughby Smith discovered the photoconductivity of the element selenium in 1873, laying the groundwork for the selenium cell which was used as a pickup in most mechanical scan systems.
In 1885, Henry Sutton in Ballarat, Australia designed what he called a telephane for transmission of images via telegraph wires, based on the Nipkow spinning disk system, selenium photocell, Nicol prisms and Kerr effect cell. Sutton's design was published internationally in 1890. An account of its use to transmit and preserve a still image was published in the Evening Star in Washington in 1896.
The first demonstration of the instantaneous transmission of images was made by a German physicist, Ernst Ruhmer, who arranged 25 selenium cells as the picture elements for a television receiver. In late 1909 he successfully demonstrated in Belgium the transmission of simple images over a telephone wire from the Palace of Justice at Brussels to the city of Liège, a distance of . This demonstration was described at the time as "the world's first working model of television apparatus". The limited number of elements meant his device was only capable of representing simple geometric shapes, and the cost was very high; at a price of £15 (US$45) per selenium cell, he estimated that a 4,000 cell system would cost £60,000 (US$180,000), and a 10,000 cell mechanism capable of reproducing "a scene or event requiring the background of a landscape" would cost £150,000 (US$450,000). Ruhmer expressed the hope that the 1910 Brussels would sponsor the construction of an advanced device with significantly more cells, as a showcase for the exposition. However, the estimated expense of £250,000 (US$750,000) proved to be too high.
The publicity generated by Ruhmer's demonstration spurred two French scientists, Georges Rignoux and A. Fournier in Paris, to announce similar research that they had been conducting. A matrix of 64 selenium cells, individually wired to a mechanical commutator, served as an electronic retina. In the receiver, a type of Kerr cell modulated the light and a series of variously angled mirrors attached to the edge of a rotating disc scanned the modulated beam onto the display screen. A separate circuit regulated synchronization. The resolution in this proof-of-concept demonstration was just sufficient to clearly transmit individual letters of the alphabet. An updated image was transmitted "several times" each second.
In 1911, Boris Rosing and his student Vladimir Zworykin created a system that used a mechanical mirror-drum scanner to transmit, in Zworykin's words, "very crude images" over wires to the "Braun tube" (cathode-ray tube or "CRT") in the receiver. Moving images were not possible because, in the scanner, "the sensitivity was not enough and the selenium cell was very laggy".
Television demonstrations
As a 23-year-old German university student, Paul Julius Gottlieb Nipkow proposed and patented the Nipkow disk in 1884. This was a spinning disk with a spiral pattern of holes in it, so each hole scanned a line of the image. Although he never built a working model of the system, Nipkow's spinning-disk "image rasterizer" was the key mechanism used in most mechanical scan systems, in both the transmitter and receiver.
Constantin Perskyi had coined the word television in a paper read to the International Electricity Congress at the International World Fair in Paris on August 24, 1900. Perskyi's paper reviewed the existing electromechanical technologies, mentioning the work of Nipkow and others.<ref>"Télévision au moyen de l'électricité", Congrès Inographs by Telegraph", The New York Times, Sunday Magazine, September 20, 1907, p. 7.</ref> However, it was the 1907 invention of the first amplifying vacuum tube, the triode, by Lee de Forest, that made the design practical.
Scottish inventor John Logie Baird in 1925 built some of the first prototype video systems, which employed the Nipkow disk. On March 25, 1925, Baird gave the first public demonstration of televised silhouette images in motion, at Selfridge's Department Store in London. Since human faces had inadequate contrast to show up on his primitive system, he televised a ventriloquist's dummy named "Stooky Bill" talking and moving, whose painted face had higher contrast. By January 26, 1926, he demonstrated the transmission of images of a face in motion by radio. This is widely regarded as being the world's first public television demonstration. Baird's system used the Nipkow disk for both scanning the image and displaying it. A brightly illuminated subject was placed in front of a spinning Nipkow disk set with lenses which swept images across a static photocell. The thallium sulphide (Thalofide) cell, developed by Theodore Case in the USA, detected the light reflected from the subject and converted it into a proportional electrical signal. This was transmitted by AM radio waves to a receiver unit, where the video signal was applied to a neon light behind a second Nipkow disk rotating synchronized with the first. The brightness of the neon lamp was varied in proportion to the brightness of each spot on the image. As each hole in the disk passed by, one scan line of the image was reproduced. Baird's disk had 30 holes, producing an image with only 30 scan lines, just enough to recognize a human face. In 1927, Baird transmitted a signal over of telephone line between London and Glasgow. In 1928, Baird's company (Baird Television Development Company/Cinema Television) broadcast the first transatlantic television signal, between London and New York, and the first shore-to-ship transmission. In 1929, he became involved in the first experimental mechanical television service in Germany. In November of the same year, Baird and Bernard Natan of Pathé established France's first television company, Télévision-Baird-Natan. In 1931, he made the first outdoor remote broadcast, of The Derby. In 1932, he demonstrated ultra-short wave television. Baird's mechanical system reached a peak of 240-lines of resolution on BBC television broadcasts in 1936 though the mechanical system did not scan the televised scene directly. Instead a 17.5 mm film was shot, rapidly developed and then scanned while the film was still wet.
An American inventor, Charles Francis Jenkins also pioneered the television. He published an article on "Motion Pictures by Wireless" in 1913, but it was not until December 1923 that he transmitted moving silhouette images for witnesses, and it was on June 13, 1925, that he publicly demonstrated synchronized transmission of silhouette pictures. In 1925 Jenkins used Nipkow disk and transmitted the silhouette image of a toy windmill in motion, over a distance of from a naval radio station in Maryland to his laboratory in Washington, D.C., using a lensed disk scanner with a 48-line resolution. He was granted the U.S. patent No. 1,544,156 (Transmitting Pictures over Wireless) on June 30, 1925 (filed March 13, 1922).
On December 25, 1925, Kenjiro Takayanagi demonstrated a television system with a 40-line resolution that employed a Nipkow disk scanner and CRT display at Hamamatsu Industrial High School in Japan. This prototype is still on display at the Takayanagi Memorial Museum in Shizuoka University, Hamamatsu Campus. By 1927, he improved the resolution to 100 lines, which was unrivaled until 1931. By 1928, he was the first to transmit human faces in half-tones. His work had an influence on the later work of Vladimir K. Zworykin. In Japan he is viewed as the man who completed the first all-electronic television. His research in creating a production model was halted by the US after Japan lost World War II.
Herbert E. Ives and Frank Gray of Bell Telephone Laboratories gave a dramatic demonstration of mechanical television on April 7, 1927. The reflected-light television system included both small and large viewing screens. The small receiver had a screen, width by height. The large receiver had a screen , width by height. Both sets were capable of reproducing reasonably accurate, monochromatic moving images. Along with the pictures, the sets also received synchronized sound. The system transmitted images over two paths: first, a copper wire link from Washington to New York City, then a radio link from Whippany, New Jersey. Comparing the two transmission methods, viewers noted no difference in quality. Subjects of the telecast included Secretary of Commerce Herbert Hoover. A flying-spot scanner beam illuminated these subjects. The scanner that produced the beam had a 50-aperture disk. The disc revolved at a rate of 18 frames per second, capturing one frame about every . (Today's systems typically transmit 30 or 60 frames per second, or one frame every respectively.) Television historian Albert Abramson underscored the significance of the Bell Labs demonstration: "It was in fact the best demonstration of a mechanical television system ever made to this time. It would be several years before any other system could even begin to compare with it in picture quality."
In 1928, General Electric launched their own experimental television station W2XB, broadcasting from the GE plant in Schenectady, New York. The station was popularly known as "WGY Television", named after the GE owned radio station WGY. The station eventually converted to an all-electronic system in the 1930s and in 1942, received a commercial license as WRGB. The station is still operating today.
Meanwhile, in the Soviet Union, Léon Theremin had been developing a mirror drum-based television, starting with 16 lines resolution in 1925, then 32 lines and eventually 64 using interlacing in 1926, and as part of his thesis on May 7, 1926, he electrically transmitted and then projected near-simultaneous moving images on a square screen. By 1927 he achieved an image of 100 lines, a resolution that was not surpassed until 1931 by RCA, with 120 lines.
Because only a limited number of holes could be made in the disks, and disks beyond a certain diameter became impractical, image resolution on mechanical television broadcasts was relatively low, ranging from about 30 lines up to 120 or so. Nevertheless, the image quality of 30-line transmissions steadily improved with technical advances, and by 1933 the UK broadcasts using the Baird system were remarkably clear. A few systems ranging into the 200-line region also went on the air.
180-lines broadcast tests were carried out by the Reichs-Rundfunk-Gesellschaft in 1935, with a transmitter in Berlin. Transmissions lasted 90 minutes a day, three days a week, with sound/visions frequencies being .
Likewise, a 180-line system that Compagnie des Compteurs (CDC) installed in Paris was tested in 1935, and a 180-line system by Peck Television Corp. started in 1935 at station VE9AK in Montreal, Quebec, Canada.
Color television
John Baird's 1928 color television experiments had inspired Goldmark's more advanced field-sequential color system. The CBS color television system invented by Peter Goldmark used such technology in 1940. In Goldmark's system, stations transmit color saturation values electronically; however, mechanical methods are also used. At the transmitting camera, a mechanical disc filters hues (colors) from reflected studio lighting. At the receiver, a synchronized disc paints the same hues over the CRT. As the viewer watches pictures through the color disc, the pictures appear in full color.
Later, simultaneous color systems superseded the CBS-Goldmark system, but mechanical color methods continued to find uses. Early color sets were very expensive: over $1,000 in the money of the time. Inexpensive adapters allowed owners of black-and-white NTSC television sets to receive color telecasts. The most prominent of these adapters is Col-R-Tel, a 1955 NTSC to field-sequential converter. This system operates at NTSC scanning rates, but uses a disc like the obsolete CBS system had. The disc converts the black-and-white set to a field-sequential set. Meanwhile, Col-R-Tel electronics recover NTSC color signals and sequence them for disc reproduction. The electronics also synchronize the disc to the NTSC system. In Col-R-Tel, the electronics provide the saturation values (chroma). These electronics cause chroma values to superimpose over brightness (luminance) changes of the picture. The disc paints the hues (color) over the picture.
A few years after Col-R-Tel, the Apollo Moon missions also adopted field-sequential techniques. The lunar color cameras all had color wheels. These Westinghouse and later RCA cameras sent field-sequential color television pictures to Earth. The Earth receiving stations included electronic equipment that converted the raw colour video signals into the NTSC standard.
Decline
The advancement of vacuum tube electronic television (including image dissectors and other camera tubes and CRTs for the reproducer) marked the beginning of the end for mechanical systems as the dominant form of television. Mechanical TV usually only produced small images. It was the main type of TV until the 1930s.
Vacuum tube television, first demonstrated in September 1927 in San Francisco by Philo Farnsworth, and then publicly by Farnsworth at the Franklin Institute in Philadelphia in 1934, was rapidly overtaking mechanical television. Farnsworth's system was first used for broadcasting in 1936, reaching 400 to more than 600 lines with fast field scan rates, along with competing systems by Philco and DuMont Laboratories. In 1939, RCA paid Farnsworth $1 million for his patents after ten years of litigation, and RCA began demonstrating all-electronic television at the 1939 World's Fair in New York City. The last mechanical television broadcasts ended in 1939 at stations run by a handful of public universities in the United States.
'Scophony' mechanical display receiver
Early Cathode-Ray Television tube displays were small in size. The 'Scophony' television receiver of 1938, an advanced television receiver that used a mechanical display, was capable of displaying a 405-line picture (compatible with the then 405-line television system used in the United Kingdom) on a display that was wide and high. A version intended for theater audiences had a wide display. It was also capable of being set up for the US 441-line television system. For 405 lines, it used a high-speed scanner running at and a low speed mirror drum running at around , in conjunction with a Jeffree cell to modulate a focused light beam from a mercury lamp. It used 39 vacuum tubes in its electronic circuits, and consumed around . Although producing impressive results and reaching the marketplace, the receiver was very expensive, costing around twice as much as a cathode-ray television. It was not a commercial success, and television transmissions in the UK were suspended for the duration of the Second World War, sealing its fate. No complete receiver survives, although some components do.
Modern applications of mechanical scanning
Since the 1970s, some amateur radio enthusiasts have experimented with mechanical systems. The early light source of a neon lamp has now been replaced with super-bright LEDs. There is some interest in creating these systems for narrow-bandwidth television, which would allow a small or large moving image to fit into a channel less than wide (modern TV systems usually have a channel about wide, 150 times larger). Also associated with this is slow-scan TV – although that typically used electronic systems utilising the P7 CRT until the 1980s and PCs thereafter. There are three known mechanical monitor forms: two fax printer-like monitors made in the 1970s, and in 2013 a small drum monitor with a coating of glow paint where the image is painted on the rotating drum with a UV laser.
Digital light processing (DLP) projectors use an array of tiny () electrostatically-actuated mirrors selectively reflecting a light source to create an image. Many low-end DLP systems also use a color wheel to provide a sequential color image, a feature that was common on many early color television systems before the shadow mask CRT provided a practical method for producing a simultaneous color image.
Another place where high-quality imagery is produced by opto-mechanics is the laser printer, where a small rotating mirror is used to deflect a modulated laser beam in one axis while the motion of the photoconductor provides the motion in the other axis. A modification of such a system using high power lasers is used in laser video projectors, with resolutions as high as 1,024 lines and each line containing over 1,500 points. Such systems produce, arguably, the best quality video images. They are used, for instance, in planetariums.
Mechanical techniques are also used in long wave infrared cameras used in military applications such as night vision for fighter pilots. These cameras use a high sensitivity infrared photo receptor (usually cooled to increase sensitivity), but instead of conventional lenses, these systems use rotating prisms to provide a 525 or 625 line standard video output. The optical parts are made from germanium, because glass is opaque at the wavelengths involved. Similar cameras have also found a role in sporting events where they are able to show (for example) where a ball has struck a bat.
Laser lighting display techniques are combined with computer emulation in the LaserMAME project. It is a vector-based system, unlike the raster displays thus-far described. Laser light reflected from computer-controlled mirrors traces out images generated by classic arcade software which is executed by a specially modified version of the MAME emulation software.
Technical aspects
Flying spot scanners
The most common method for creating the video signal was the "flying spot scanner", developed as a remedy for the low sensitivity that photoelectric cells had at the time. Instead of a television camera that took pictures, a flying spot scanner projected a bright spot of light that scanned rapidly across the subject scene in a raster pattern, in a darkened studio. The light reflected from the subject was picked up by banks of photoelectric cells and amplified to become the video signal.
In the scanner the narrow light beam was produced by an arc lamp shining through the holes in a spinning Nipkow disk. Each sweep of the spot across the scene produced a "scan line" of the picture. A single "frame" of the picture was typically made up of 24, 48, or 60 scan lines. The scene was typically scanned 15 or 20 times per second, producing 15 or 20 video frames per second. The varying brightness of the point where the spot fell reflected varying amounts of light, which was converted to a proportionally varying electronic signal by the photoelectric cells. To achieve adequate sensitivity, instead of a single cell, a number of photoelectric cells were used. Like mechanical television itself, flying spot technology grew out of phototelegraphy (facsimile). This scanning method began in the 19th century.
The flying spot method has two disadvantages:
Actors must perform in near darkness
Flying spot cameras tend to work unreliably outdoors in daylight
In 1928, Ray Kell from the United States' General Electric proved that flying spot scanners could work outdoors. The scanning light source must be brighter than other incident illumination.
Kell was the engineer who ran a 24-line camera that telecast pictures of New York governor Al Smith. Smith was accepting the Democratic nomination for presidency. As Smith stood outside the capital in Albany, Kell managed to send usable pictures to his associate Bedford at station WGY, which was broadcasting Smith's speech. The rehearsal went well, but then the real event began. The newsreel cameramen switched on their floodlights.
Unfortunately for Kell, his scanner only had a lamp inside it. The floodlights threw much more light on Governor Smith. These floods simply overwhelmed Kell's imaging photocells. In fact, the floods made the unscanned part of the image as bright as the scanned part. Kell's photocells couldn't discriminate reflections off Smith (from the AC scanning beam) from the flat, DC light from the floodlamps.
The effect is very similar to extreme overexposure in a still camera: The scene disappears, and the camera records a flat, bright light. If used in favorable conditions, however, the picture comes out correctly. Similarly, Kell proved that outdoors in favorable conditions, his scanner worked.
The BBC television service used the flying spot method until 1935, and German television used flying spot methods as late as 1938. However, flying spot techniques remained in use in many applications after the demise of mechanical television. The German inventor Manfred von Ardenne designed a flying spot scanner with a CRT as the light source, and CRT-based flying spot scanners became a common technique for telecine. In the 1950s, DuMont marketed Vitascan, an entire flying-spot color studio system. Laser scanners continue to use a flying spot approach.
Larger videos
A few mechanical TV systems could produce images several feet or meters wide and of comparable quality to the CRT televisions that were to follow. CRT technology at that time was limited to small, low-brightness screens. One such system was developed by Ulises Armand Sanabria in Chicago. By 1934, Sanabria demonstrated a projection system which had a image.
Perhaps the best mechanical televisions of the 1930s used the Scophony system, which could produce images of more than 400 lines and display them on screens at least in size (at least a few models of this type were actually produced).
The Scophony system used multiple drums rotating at fairly high speed to create the images. One using a 441-line American standard of the day had a small drum rotating at (a second slower drum moved at just a few hundred rpm).
Aspect ratios
Some mechanical equipment scanned lines vertically rather than horizontally, as in modern TVs. An example of this method is the Baird 30-line system. Baird's British system created a picture in the shape of a very narrow, vertical rectangle.
This shape created a "portrait" image, instead of the "landscape" orientationthese terms coming from the concepts of portrait and landscape in artthat is common today. The position of a framing mask before the Nipkow disk determines the scan line orientation. Placement of the framing mask at the left or right side of the disk gives vertical scan lines. Placement at the top or bottom of the disk gives horizontal scan lines.
Baird's earliest television images had very low definition. These images could only show one person clearly. For this reason, a vertical "portrait" image made more sense to Baird than a horizontal "landscape" image. Baird chose a shape three units wide by seven high. This shape is only about half as wide as a traditional portrait and close in proportion to a typical doorway.
Instead of entertainment television, Baird might have had point-to-point communication in mind. Another television system followed that reasoning. The 1927 system developed by Herbert E. Ives at AT&T's Bell Laboratories was a large-screen television system and the most advanced television of its day. The Ives 50-line system also produced a vertical "portrait" picture. Since AT&T intended to use television for telephony, the vertical shape was logical: phone calls are usually conversations between just two people. A picturephone system would depict one person on each side of the line.
Meanwhile, in the US, Germany and elsewhere, other inventors planned to use television for entertainment purposes. These inventors began with square or "landscape" pictures. (For example, the television systems of Ernst Alexanderson, Frank Conrad, Charles Francis Jenkins, William Peck and Ulises Armand Sanabria.) These inventors realized that television is about relationships between people. From the very beginning, these inventors allowed picture space for two-shots. Soon, images increased to 60 lines or more. The camera could easily photograph several people at once. Then even Baird switched his picture mask to a horizontal image. Baird's "zone television" is an early example of rethinking his extremely narrow screen format. For entertainment and most other purposes, even today, landscape remains the more practical shape.
Recording
In the days of commercial mechanical television transmissions, a system of recording images (but not sound) was developed, using a modified gramophone recorder. Marketed as "Phonovision", this system, which was never fully perfected, proved to be complicated to use as well as quite expensive, yet managed to preserve a number of early broadcast images that would otherwise have been lost. Scottish computer engineer Donald F. McLean has painstakingly reconstructed the analogue playback technology required to view these recordings, and has given lectures and presentations on his collection of mechanical television recordings made between 1925 and 1933.
Among the discs in Dr. McLean's collection are a number of test recordings made by television pioneer John Logie Baird himself. One disc, dated "28th March 1928" and marked with the title "Miss Pounsford", shows several minutes of a woman's face in what appears to be very animated conversation. In 1993, the woman was identified by relatives as Mabel Pounsford, and her brief appearance on the disc is one of the earliest known television video recordings of a human.
Bibliography
Beyer, Rick, The Greatest Stories Never Told : 100 tales from history to astonish, bewilder, & stupefy, A&E Television Networks, 2003,
Huurdeman, Anton A., The worldwide history of telecommunications, Wiley-IEEE, 2003,
Sarkar, Tapan K. et al., History of wireless, John Wiley and Sons, 2006,
Simonis, Doris, Inventors and Inventions'', Marshall Cavendish, 2007, – via Google Books
| Technology | Broadcasting | null |
1362038 | https://en.wikipedia.org/wiki/Quaternary%20ammonium%20cation | Quaternary ammonium cation | In organic chemistry, quaternary ammonium cations, also known as quats, are positively-charged polyatomic ions of the structure , where R is an alkyl group, an aryl group or organyl group. Unlike the ammonium ion () and the primary, secondary, or tertiary ammonium cations, the quaternary ammonium cations are permanently charged, independent of the pH of their solution. Quaternary ammonium salts or quaternary ammonium compounds (called quaternary amines in oilfield parlance) are salts of quaternary ammonium cations. Polyquats are a variety of engineered polymer forms which provide multiple quat molecules within a larger molecule.
Quats are used in consumer applications including as antimicrobials (such as detergents and disinfectants), fabric softeners, and hair conditioners. As an antimicrobial, they are able to inactivate enveloped viruses (such as SARS-CoV-2). Quats tend to be gentler on surfaces than bleach-based disinfectants, and are generally fabric-safe.
Synthesis
Quaternary ammonium compounds are prepared by the alkylation of tertiary amine. Industrial production of commodity quat salts usually involves hydrogenation of fatty nitriles, which can generate primary or secondary amines. These amines are then treated with methyl chloride.
The quaternization of alkyl amines by alkyl halides is widely documented. In older literature this is often called a Menshutkin reaction, however modern chemists usually refer to it simply as quaternization. The reaction can be used to produce a compound with unequal alkyl chain lengths; for example when making cationic surfactants one of the alkyl groups on the amine is typically longer than the others. A typical synthesis is for benzalkonium chloride from a long-chain alkyldimethylamine and benzyl chloride:
Reactions
Quaternary ammonium cations are unreactive toward even strong electrophiles, oxidants, and acids. They also are stable toward most nucleophiles. The latter is indicated by the stability of the hydroxide salts such as tetramethylammonium hydroxide and tetrabutylammonium hydroxide even at elevated temperatures. The halflife of Me4NOH in 6M NaOH at 160 °C is >61 h.
Because of their resilience, many unusual anions have been isolated as the quaternary ammonium salts. Examples include tetramethylammonium pentafluoroxenate, containing the highly reactive pentafluoroxenate () ion. Permanganate can be solubilized in organic solvents, when deployed as its salt.
With exceptionally strong bases, quat cations degrade. They undergo Sommelet–Hauser rearrangement and Stevens rearrangement, as well as dealkylation under harsh conditions or in presence of strong nucleophiles, like thiolates. Quaternary ammonium cations containing N−C−C−H units can also undergo the Hofmann elimination and Emde degradation.
Examples
Tetramethylammonium ion: , also denoted (Me = methyl group)
Tetraethylammonium ion: , also denoted (Et = ethyl group)
Tetrapropylammonium ion: , also denoted (Pr = propyl group)
Tetrabutylammonium ion: , also denoted (Bu = butyl group)
Applications
Quaternary ammonium salts are used as disinfectants, surfactants, fabric softeners, and as antistatic agents (e.g. in shampoos). In liquid fabric softeners, the chloride salts are often used. In dryer anticling strips, the sulfate salts are often used. Older aluminium electrolytic capacitors and spermicidal jellies also contain quaternary ammonium salts. Quats are also used in contraception formulations, veterinary products, diagnostic testing, vaccine production, and nasal formulations.
Concerns have been raised about the level of understanding of safety profile of quat disinfectants on people. As of August 2020, half of disinfectants the United States Environmental Protection Agency suggested as effective against COVID-19 contained one of the quats, and often a quat as the sole ingredient. Salmonella and E. coli O157:H7 exposed to quats have developed cross resistance to antibiotics. A subject of concern is the potential effect of increased use of quats related to COVID-19 pandemic on antibiotic resistance in a larger microbial community in nature and engineered environment.
Medicines
Quaternary ammonium compounds have antimicrobial activity. Quaternary ammonium compounds, especially those containing long alkyl chains, are used as antimicrobials and disinfectants. Examples are benzalkonium chloride, benzethonium chloride, methylbenzethonium chloride, cetalkonium chloride, cetylpyridinium chloride, cetrimonium, cetrimide, dofanium chloride, tetraethylammonium bromide, didecyldimethylammonium chloride and domiphen bromide. Also good against fungi, amoebas, and enveloped viruses (such as SARS-CoV-2), most quaternary ammonium compounds are believed to act by disrupting the cell membrane or viral envelope. (Some QACs, such as dequalinium and similar bis-QACs, show evidence of a different mode of action.)
Quaternary ammonium compounds are lethal to a wide variety of organisms except endospores and non-enveloped viruses, both having no accessible membrane coat to attack. It is possible to solve the endospore problem by adding chemicals which force them to germinate. They have reduced efficacy against gram-negative bacteria, mycobacteria, and bacteria in biofilms due to them having additional layers that need to be penetrated or disrupted. Some bacteria such as MRSA have acquired resistance genes, qacA/B and qacC/D, that pump the cation out of the cell.
Phase transfer catalysts
In organic chemistry, quaternary ammonium salts are employed as phase transfer catalysts (PTCs). Such catalysts accelerate reactions between reagents dissolved in immiscible solvents. The highly reactive reagent dichlorocarbene is generated via PTC by reaction of chloroform and aqueous sodium hydroxide.
Fabric softeners and hair conditioners
In the 1950s, distearyldimethylammonium chloride (DHTDMAC), was introduced as a fabric softener. This compound was discontinued because the cation biodegrades too slowly. Contemporary fabric softeners are based on salts of quaternary ammonium cations where the fatty acid is linked to the quaternary center via ester linkages; these are commonly referred to as betaine-esters or ester-quats and are susceptible to degradation, e.g., by hydrolysis. Characteristically, the cations contain one or two long alkyl chains derived from fatty acids linked to an ethoxylated ammonium salt. Other cationic compounds can be derived from imidazolium, guanidinium, substituted amine salts, or quaternary alkoxy ammonium salts.
The antistatic qualities that make quaternary ammonium salts useful as fabric softeners also make them useful in hair conditioners and shampoos. The idea was pioneered by Henkel with a 1984 patent. Examples include cetrimonium chloride and behentrimonium chloride.
Plant growth retardants
Cycocel (chlormequat chloride) reduces plant height by inhibiting the production of gibberellins, the primary plant hormones responsible for cell elongation. Therefore, their effects are primarily on stem, petiole, and flower stalk tissues. Lesser effects are seen in reductions of leaf expansion, resulting in thicker leaves with darker green color.
Natural occurrence
Several quaternary ammonium derivatives exist in nature. Prominent examples include glycine betaine, choline, carnitine, butyrobetaine, homarine, and trigonelline. Glycine betaine, an osmolyte, stabilizes osmotic pressure in cells.
Choline is a precursor for the neurotransmitter acetylcholine. Choline is also a constituent of lecithin, which is present in many plants and animal organs. It is found in phospholipids. For example, phosphatidylcholines, a major component of biological membranes, are a member of the lecithin group of fatty substances in animal and plant tissues.
Carnitine participates in the beta-oxidation of fatty acids.
Health effects
Quaternary ammonium compounds can display a range of health effects, amongst which are mild skin and respiratory irritation up to severe caustic burns on skin and the gastrointestinal wall (depending on concentration), gastrointestinal symptoms (e.g., nausea and vomiting), coma, convulsions, hypotension and death.
They are thought to be the chemical group responsible for anaphylactic reactions that occur with use of neuromuscular blocking drugs during general anaesthesia in surgery. Quaternium-15 is the single most often found cause of allergic contact dermatitis of the hands (16.5% in 959 cases).
Possible reproductive effects in laboratory animals
Quaternary ammonium-based disinfectants (Virex and Quatricide) were tentatively identified as the most probable cause of jumps in birth defects and fertility problems in caged lab mice. The quat ingredients in the disinfectants include alkyl dimethyl benzyl ammonium chloride (ADBAC) and didecyl dimethyl ammonium chloride (DDAC). A similar link was tentatively identified in nurses. The studies contradict earlier toxicology data reviewed by the U.S. Environmental Protection Agency (U.S. EPA) and the EU Commission.
Quantification
The quantification of quaternary ammonium compounds can be challenging. Some methods include precipitation of solid salts with tetraphenylborate. Another method, an Epton titration, involves partitioning between water-chloroform in the presence of an anionic dye. Individual cations are detectable by ESI-MS and NMR spectroscopy.
| Physical sciences | Salts and ions: General | Chemistry |
1362114 | https://en.wikipedia.org/wiki/Aluminium%20chloride | Aluminium chloride | Aluminium chloride, also known as aluminium trichloride, is an inorganic compound with the formula . It forms a hexahydrate with the formula , containing six water molecules of hydration. Both the anhydrous form and the hexahydrate are colourless crystals, but samples are often contaminated with iron(III) chloride, giving them a yellow colour.
The anhydrous form is commercially important. It has a low melting and boiling point. It is mainly produced and consumed in the production of aluminium, but large amounts are also used in other areas of the chemical industry. The compound is often cited as a Lewis acid. It is an example of an inorganic compound that reversibly changes from a polymer to a monomer at mild temperature.
History
The salt was known in the 18th century as muriate of alumina, marine alum, argillaceous marine salt, muriated clay. It was first chemically studied in the 1830s.
Structure
Anhydrous
adopts three structures, depending on the temperature and the state (solid, liquid, gas). Solid has a sheet-like layered structure with cubic close-packed chloride ions. In this framework, the Al centres exhibit octahedral coordination geometry. Yttrium(III) chloride adopts the same structure, as do a range of other compounds. When aluminium trichloride is in its melted state, it exists as the dimer , with tetracoordinate aluminium. This change in structure is related to the lower density of the liquid phase (1.78 g/cm3) versus solid aluminium trichloride (2.48 g/cm3). dimers are also found in the vapour phase. At higher temperatures, the dimers dissociate into trigonal planar monomer, which is structurally analogous to . The melt conducts electricity poorly, unlike more ionic halides such as sodium chloride.
Aluminium chloride monomer belongs to the point group D3h in its monomeric form and D2h in its dimeric form.
Hexahydrate
The hexahydrate consists of octahedral cation centers and chloride anions () as counterions. Hydrogen bonds link the cation and anions.
The hydrated form of aluminium chloride has an octahedral molecular geometry, with the central aluminium ion surrounded by six water ligand molecules. Being coordinatively saturated, the hydrate is of little value as a catalyst in Friedel-Crafts alkylation and related reactions.
Uses
Alkylation and acylation of arenes
is a common Lewis-acid catalyst for Friedel-Crafts reactions, both acylations and alkylations. Important products are detergents and ethylbenzene. These types of reactions are the major use for aluminium chloride, for example, in the preparation of anthraquinone (used in the dyestuffs industry) from benzene and phosgene. In the general Friedel-Crafts reaction, an acyl chloride or alkyl halide reacts with an aromatic system as shown:
The alkylation reaction is more widely used than the acylation reaction, although its practice is more technically demanding. For both reactions, the aluminium chloride, as well as other materials and the equipment, should be dry, although a trace of moisture is necessary for the reaction to proceed. Detailed procedures are available for alkylation and acylation of arenes.
A general problem with the Friedel-Crafts reaction is that the aluminium chloride catalyst sometimes is required in full stoichiometric quantities, because it complexes strongly with the products. This complication sometimes generates a large amount of corrosive waste. For these and similar reasons, the use of aluminium chloride has often been displaced by zeolites.
Aluminium chloride can also be used to introduce aldehyde groups onto aromatic rings, for example via the Gattermann-Koch reaction which uses carbon monoxide, hydrogen chloride and a copper(I) chloride co-catalyst.
Other applications in organic and organometallic synthesis
Aluminium chloride finds a wide variety of other applications in organic chemistry. For example, it can catalyse the ene reaction, such as the addition of 3-buten-2-one (methyl vinyl ketone) to carvone:
It is used to induce a variety of hydrocarbon couplings and rearrangements.
Aluminium chloride combined with aluminium in the presence of an arene can be used to synthesize bis(arene) metal complexes, e.g. bis(benzene)chromium, from certain metal halides via the Fischer–Hafner synthesis. Dichlorophenylphosphine is prepared by reaction of benzene and phosphorus trichloride catalyzed by aluminium chloride.
Medical
Topical aluminum chloride hexahydrate is used for the treatment of hyperhidrosis (excessive sweating).
Reactions
Anhydrous aluminium chloride is a powerful Lewis acid, capable of forming Lewis acid-base adducts with even weak Lewis bases such as benzophenone and mesitylene. It forms tetrachloroaluminate () in the presence of chloride ions.
Aluminium chloride reacts with calcium and magnesium hydrides in tetrahydrofuran forming tetrahydroaluminates.
Reactions with water
Anhydrous aluminium chloride is hygroscopic, having a very pronounced affinity for water. It fumes in moist air and hisses when mixed with liquid water as the Cl− ligands are displaced with H2O molecules to form the hexahydrate . The anhydrous phase cannot be regained on heating the hexahydrate. Instead HCl is lost leaving aluminium hydroxide or alumina (aluminium oxide):
Like metal aquo complexes, aqueous is acidic owing to the ionization of the aquo ligands:
Aqueous solutions behave similarly to other aluminium salts containing hydrated ions, giving a gelatinous precipitate of aluminium hydroxide upon reaction with dilute sodium hydroxide:
Synthesis
Aluminium chloride is manufactured on a large scale by the exothermic reaction of aluminium metal with chlorine or hydrogen chloride at temperatures between .
Aluminium chloride may be formed via a single displacement reaction between copper(II) chloride and aluminium.
In the US in 1993, approximately 21,000 tons were produced, not counting the amounts consumed in the production of aluminium.
Hydrated aluminium trichloride is prepared by dissolving aluminium oxides in hydrochloric acid. Metallic aluminium also readily dissolves in hydrochloric acid ─ releasing hydrogen gas and generating considerable heat. Heating this solid does not produce anhydrous aluminium trichloride, the hexahydrate decomposes to aluminium hydroxide when heated:
Aluminium also forms a lower chloride, aluminium(I) chloride (AlCl), but this is very unstable and only known in the vapour phase.
Natural occurrence
Anhydrous aluminium chloride is not found as a mineral. The hexahydrate, however, is known as the rare mineral chloraluminite. A more complex, basic and hydrated aluminium chloride mineral is cadwaladerite.
Safety
Anhydrous reacts vigorously with bases, so suitable precautions are required.
It can cause irritation to the eyes, skin, and the respiratory system if inhaled or on contact.
| Physical sciences | Halide salts | Chemistry |
1362656 | https://en.wikipedia.org/wiki/Eugenol | Eugenol | Eugenol is an allyl chain-substituted guaiacol, a member of the allylbenzene class of chemical compounds. It is a colorless to pale yellow, aromatic oily liquid extracted from certain essential oils especially from clove, nutmeg, cinnamon, basil and bay leaf. It is present in concentrations of 80–90% in clove bud oil and at 82–88% in clove leaf oil. Eugenol has a pleasant, spicy, clove-like scent. The name is derived from Eugenia caryophyllata, the former Linnean nomenclature term for cloves. The currently accepted name is Syzygium aromaticum.
Biosynthesis
The biosynthesis of eugenol begins with the amino acid tyrosine. L-tyrosine is converted to p-coumaric acid by the enzyme tyrosine ammonia lyase (TAL). From here, p-coumaric acid is converted to caffeic acid by p-coumarate 3-hydroxylase using oxygen and NADPH. S-Adenosyl methionine (SAM) is then used to methylate caffeic acid, forming ferulic acid, which is in turn converted to feruloyl-CoA by the enzyme 4-hydroxycinnamoyl-CoA ligase (4CL). Next, feruloyl-CoA is reduced to coniferaldehyde by cinnamoyl-CoA reductase (CCR). Coniferaldeyhyde is then further reduced to coniferyl alcohol by cinnamyl-alcohol dehydrogenase (CAD) or sinapyl-alcohol dehydrogenase (SAD). Coniferyl alcohol is then converted to an ester in the presence of the substrate CH3COSCoA, forming coniferyl acetate. Finally, coniferyl acetate is converted to eugenol via the enzyme eugenol synthase 1 and the use of NADPH.
Eugenol is a metabolite of Caleicine, the active compound found in Calea Ternifolia, and is thought to cause the sedative and hallucinogenic state Calea Ternifolia can induce.
Pharmacology
Eugenol and thymol possess general anesthetic properties. Like many other anesthetic agents, these 2-alkyl(oxy)phenols act as positive allosteric modulators of the GABAA receptor. Although eugenol and thymol are too toxic and not potent enough to be used clinically, these findings led to the development of 2-substituted phenol anesthetic drugs, including propanidid (later withdrawn) and the widely used propofol. Eugenol and the structurally similar myristicin, have the common property of inhibiting MAO-A and MAO-B in vitro.
Eugenol acts on the NMDA receptors as an antagonist, as well as the histamine receptors as an antagonist, but the exact specifics of this are unknown.
In humans, complete excretion occurs within 24 hour and metabolites are mostly conjugates of eugenol.
Uses
Humans
Eugenol is used as a flavor or aroma ingredient in teas, meats, cakes, perfumes, cosmetics, flavorings, and essential oils. It is also used as a local antiseptic and anaesthetic. Eugenol can be combined with zinc oxide to form zinc oxide eugenol which has restorative and prosthodontic applications in dentistry. For persons with a dry socket as a complication of tooth extraction, packing the dry socket with a eugenol-zinc oxide paste on iodoform gauze is effective for reducing acute pain. Eugenol-zinc oxide paste is also used for root canal sealing.
Insects and fish
It is attractive to males of various species of orchid bees, which apparently gather the chemical to synthesize pheromones; it is commonly used as bait to attract and collect these bees for study. It also attracts female cucumber beetles.
Eugenol and isoeugenol, which both are floral volatile scent compounds, are catalyzed by a single type of enzyme in the genus Gymnadenia and the gene encoding for this enzyme is the first functionally characterized gene in these species. Eugenol is an ingredient in some insecticides.
Clove oil is common as an anesthetic for use on aquarium fish as well as on wild fish when sampled for research and management purposes. Where readily available, it presents a humane method to euthanize sick and diseased fish either by direct overdose or to induce sleep before an overdose of eugenol.
Other
Eugenol is an ingredient in some fungicides and weed control products used in agricultural practices in the European Union. It is used in hundreds of household products, such as pesticides, pet care, laundry, cleaning, and paper or vehicle products.
Toxicity
Lesser side effects of eugenol toxicity that may not be considered a full overdose:
Sedation, dizziness, hallucinations, mild respiratory depression, nausea, and muscle spasms.
Taken orally in high doses for chronic periods, eugenol may cause liver toxicity. An overdose is possible, causing a wide range of symptoms, such as hematuria (blood in urine), convulsions, diarrhea, delirium, unconsciousness, heavy respiratory depression, tachycardia (rapid heart rate), or acute kidney injury. N-acetylcysteine may be used to treat people with eugenol or clove oil overdose.
As an allergenic
Eugenol is subject to restrictions on its use in perfumery, as some people may become sensitised to it, however, the degree to which eugenol can cause an allergic reaction in humans is disputed.
Eugenol is a component of balsam of Peru, to which some people are allergic. When eugenol is used in dental preparations such as surgical pastes, dental packing, and dental cement, it may cause contact stomatitis and allergic cheilitis. The allergy can be discovered via a patch test.
Natural occurrence
Eugenol naturally occurs in numerous plants, including the following:
Cloves (Syzygium aromaticum)
Wormwood
Cinnamon
Cinnamomum tamala
Nutmeg (Myristica fragrans)
Ocimum basilicum (sweet basil)
Ocimum gratissimum (African basil)
Ocimum tenuiflorum (syn. Ocimum sanctum, tulsi or holy basil)
Japanese star anise
Lemon balm
Dill
Pimenta dioica (Allspice)
Vanilla
Bay laurel
Celery
Ginger
Wood avens
| Physical sciences | Phenylpropanoids | Chemistry |
1362795 | https://en.wikipedia.org/wiki/4-manifold | 4-manifold | In mathematics, a 4-manifold is a 4-dimensional topological manifold. A smooth 4-manifold is a 4-manifold with a smooth structure. In dimension four, in marked contrast with lower dimensions, topological and smooth manifolds are quite different. There exist some topological 4-manifolds which admit no smooth structure, and even if there exists a smooth structure, it need not be unique (i.e. there are smooth 4-manifolds which are homeomorphic but not diffeomorphic).
4-manifolds are important in physics because in general relativity, spacetime is modeled as a pseudo-Riemannian 4-manifold.
Topological 4-manifolds
The homotopy type of a simply connected compact 4-manifold only depends on the intersection form on the middle dimensional homology. A famous theorem of implies that the homeomorphism type of the manifold only depends on this intersection form, and on a invariant called the Kirby–Siebenmann invariant, and moreover that every combination of unimodular form and Kirby–Siebenmann invariant can arise, except that if the form is even, then the Kirby–Siebenmann invariant must be the signature/8 (mod 2).
Examples:
In the special case when the form is 0, this implies the 4-dimensional topological Poincaré conjecture.
If the form is the E8 lattice, this gives a manifold called the E8 manifold, a manifold not homeomorphic to any simplicial complex.
If the form is , there are two manifolds depending on the Kirby–Siebenmann invariant: one is 2-dimensional complex projective space, and the other is a fake projective space, with the same homotopy type but not homeomorphic (and with no smooth structure).
When the rank of the form is greater than about 28, the number of positive definite unimodular forms starts to increase extremely rapidly with the rank, so there are huge numbers of corresponding simply connected topological 4-manifolds (most of which seem to be of almost no interest).
Freedman's classification can be extended to some cases when the fundamental group is not too complicated; for example, when it is , there is a classification similar to the one above using Hermitian forms over the group ring of . If the fundamental group is too large (for example, a free group on 2 generators), then Freedman's techniques seem to fail and very little is known about such manifolds.
For any finitely presented group it is easy to construct a (smooth) compact 4-manifold with it as its fundamental group. (More specifically, for any finitely presented group, one constructs a manifold with the given fundamental group, such that two manifolds in this family are homeomorphic if and only if the fundamental groups are isomorphic.) As there can be no algorithm to tell whether two finitely presented groups are isomorphic (even if one is known to be trivial), there can be no algorithm to tell if two 4-manifolds have the same fundamental group. This is one reason why much of the work on 4-manifolds just considers the simply connected case: the general case of many problems is already known to be intractable.
Smooth 4-manifolds
For manifolds of dimension at most 6, any piecewise linear (PL) structure can be smoothed in an essentially unique way, so in particular the theory of 4 dimensional PL manifolds is much the same as the theory of 4 dimensional smooth manifolds.
A major open problem in the theory of smooth 4-manifolds is to classify the simply connected compact ones.
As the topological ones are known, this breaks up into two parts:
Which topological manifolds are smoothable?
Classify the different smooth structures on a smoothable manifold.
There is an almost complete answer to the first problem asking which simply connected compact 4-manifolds have smooth structures.
First, the Kirby–Siebenmann class must vanish.
If the intersection form is definite Donaldson's theorem gives a complete answer: there is a smooth structure if and only if the form is diagonalizable.
If the form is indefinite and odd there is a smooth structure.
If the form is indefinite and even we may as well assume that it is of nonpositive signature by changing orientations if necessary, in which case it is isomorphic to a sum of m copies of II1,1 and 2n copies of E8(−1) for some m and n. If m ≥ 3n (so that the dimension is at least 11/8 times the |signature|) then there is a smooth structure, given by taking a connected sum of n K3 surfaces and m − 3n copies of S2×S2. If m ≤ 2n (so the dimension is at most 10/8 times the |signature|) then Furuta proved that no smooth structure exists . This leaves a small gap between 10/8 and 11/8 where the answer is mostly unknown. (The smallest case not covered above has n=2 and m=5, but this has also been ruled out, so the smallest lattice for which the answer is not currently known is the lattice II7,55 of rank 62 with n=3 and m=7. See for recent (as of 2019) progress in this area.) The "11/8 conjecture" states that smooth structures do not exist if the dimension is less than 11/8 times the |signature|.
In contrast, very little is known about the second question of classifying the smooth structures on a smoothable 4-manifold; in fact, there is not a single smoothable 4-manifold where the answer is fully known. Donaldson showed that there are some simply connected compact 4-manifolds, such as Dolgachev surfaces, with a countably infinite number of different smooth structures. There are an uncountable number of different smooth structures on R4; see exotic R4.
Fintushel and Stern showed how to use surgery to construct large numbers of different smooth structures (indexed by arbitrary integral polynomials) on many different manifolds, using Seiberg–Witten invariants to show that the smooth structures are different. Their results suggest that any classification of simply connected smooth 4-manifolds will be very complicated. There are currently no plausible conjectures about what this classification might look like. (Some early conjectures that all simply connected smooth 4-manifolds might be connected sums of algebraic surfaces, or symplectic manifolds, possibly with orientations reversed, have been disproved.)
Special phenomena in 4 dimensions
There are several fundamental theorems about manifolds that can be proved by low-dimensional methods in dimensions at most 3, and by completely different high-dimensional methods in dimension at least 5, but which are false in dimension 4. Here are some examples:
In dimensions other than 4, the Kirby–Siebenmann invariant provides the obstruction to the existence of a PL structure; in other words a compact topological manifold has a PL structure if and only if its Kirby–Siebenmann invariant in H4(M,Z/2Z) vanishes. In dimension 3 and lower, every topological manifold admits an essentially unique PL structure. In dimension 4 there are many examples with vanishing Kirby–Siebenmann invariant but no PL structure.
In any dimension other than 4, a compact topological manifold has only a finite number of essentially distinct PL or smooth structures. In dimension 4, compact manifolds can have a countably-infinite number of non-diffeomorphic smooth structures.
Four is the only dimension n for which Rn can have an exotic smooth structure. R4 has an uncountable number of exotic smooth structures; see exotic R4.
The solution to the smooth Poincaré conjecture is known in all dimensions other than 4 (it is usually false in dimensions at least 7; see exotic sphere). The Poincaré conjecture for PL manifolds has been proved for all dimensions other than 4. In 4 dimensions, the PL Poincaré conjecture is equivalent to the smooth Poincaré conjecture, and its truth is unknown.
The smooth h-cobordism theorem holds for cobordisms provided that neither the cobordism nor its boundary has dimension 4. It can fail if the boundary of the cobordism has dimension 4 (as shown by Donaldson). If the cobordism has dimension 4, then it is unknown whether the h-cobordism theorem holds.
A topological manifold of dimension not equal to 4 has a handlebody decomposition. Manifolds of dimension 4 have a handlebody decomposition if and only if they are smoothable.
There are compact 4-dimensional topological manifolds that are not homeomorphic to any simplicial complex. Ciprian Manolescu showed that there are topological manifolds in each dimension greater than or equal to 5, that are not homeomorphic to a simplicial complex.
Failure of the Whitney trick in dimension 4
According to Frank Quinn, "Two n-dimensional submanifolds of a manifold of dimension 2n will usually intersect themselves and each other in isolated points. The "Whitney trick" uses an isotopy across an embedded 2-disk to simplify these intersections. Roughly speaking this reduces the study of n-dimensional embeddings to embeddings of 2-disks. But this is not a reduction when the dimension is 4: the 2-disks themselves are middle-dimensional, so trying to embed them encounters exactly the same problems they are supposed to solve. This is the phenomenon that separates dimension 4 from others."
| Mathematics | Topology | null |
1363291 | https://en.wikipedia.org/wiki/Medical%20device | Medical device | A medical device is any device intended to be used for medical purposes. Significant potential for hazards are inherent when using a device for medical purposes and thus medical devices must be proved safe and effective with reasonable assurance before regulating governments allow marketing of the device in their country. As a general rule, as the associated risk of the device increases the amount of testing required to establish safety and efficacy also increases. Further, as associated risk increases the potential benefit to the patient must also increase.
Discovery of what would be considered a medical device by modern standards dates as far back as in Baluchistan where Neolithic dentists used flint-tipped drills and bowstrings. Study of archeology and Roman medical literature also indicate that many types of medical devices were in widespread use during the time of ancient Rome. In the United States it was not until the Federal Food, Drug, and Cosmetic Act (FD&C Act) in 1938 that medical devices were regulated. Later in 1976, the Medical Device Amendments to the FD&C Act established medical device regulation and oversight as we know it today in the United States. Medical device regulation in Europe as we know it today came into effect in 1993 by what is collectively known as the Medical Device Directive (MDD). On May 26, 2017, the Medical Device Regulation (MDR) replaced the MDD.
Medical devices vary in both their intended use and indications for use. Examples range from simple, low-risk devices such as tongue depressors, medical thermometers, disposable gloves, and bedpans to complex, high-risk devices that are implanted and sustain life. One example of high-risk devices are those with embedded software such as pacemakers, and which assist in the conduct of medical testing, implants, and prostheses. The design of medical devices constitutes a major segment of the field of biomedical engineering.
The global medical device market was estimated to be between $220 and US$250 billion in 2013. The United States controls ≈40% of the global market followed by Europe (25%), Japan (15%), and the rest of the world (20%). Although collectively Europe has a larger share, Japan has the second largest country market share. The largest market shares in Europe (in order of market share size) belong to Germany, Italy, France, and the United Kingdom. The rest of the world comprises regions like (in no particular order) Australia, Canada, China, India, and Iran. This article discusses what constitutes a medical device in these different regions and throughout the article these regions will be discussed in order of their global market share.
Definition
A global definition for medical device is difficult to establish because there are numerous regulatory bodies worldwide overseeing the marketing of medical devices. Although these bodies often collaborate and discuss the definition in general, there are subtle differences in wording that prevent a global harmonization of the definition of a medical device, thus the appropriate definition of a medical device depends on the region. Often a portion of the definition of a medical device is intended to differentiate between medical devices and drugs, as the regulatory requirements of the two are different. Definitions also often recognize In vitro diagnostics as a subclass of medical devices and establish accessories as medical devices.
Definitions by region
United States (Food and Drug Administration)
Section 201(h) of the Federal Food Drug & Cosmetic (FD&C) Act defines a device as an "instrument, apparatus, implement, machine, contrivance, implant, in vitro reagent, or other similar or related article, including a component part, or accessory which is:
recognized in the official National Formulary, or the United States Pharmacopoeia, or any supplement to them
Intended for use in the diagnosis of disease or other conditions, or in the cure, mitigation, treatment, or prevention of disease, in man or other animals, or
Intended to affect the structure or any function of the body of man or other animals, and
which does not achieve its primary intended purposes through chemical action within or on the body of man or other animals and which is not dependent upon being metabolized for the achievement of its primary intended purposes. The term 'device' does not include software functions excluded pursuant to section 520(o)."
European Union
According to Article 1 of Council Directive 93/42/EEC, 'medical device' means any "instrument, apparatus, appliance, software, material or other article, whether used alone or in combination, including the software intended by its manufacturer to be used specifically for diagnostic and/or therapeutic purposes and necessary for its proper application, intended by the manufacturer to be used for human beings for the purpose of:
diagnosis, prevention, monitoring, treatment or alleviation of disease,
diagnosis, monitoring, treatment, alleviation of or compensation for an injury or handicap,
investigation, replacement or modification of the anatomy or of a physiological process,
control of conception,
and which does not achieve its principal intended action in or on the human body by pharmacological, immunological or metabolic means, but which may be assisted in its function by such means;"
EU Legal framework
Based on the New Approach, rules that relate to safety and performance of medical devices were harmonised in the EU in the 1990s. The New Approach, defined in a European Council Resolution of May 1985, represents an innovative way of technical harmonisation. It aims to remove technical barriers to trade and dispel the consequent uncertainty for economic operators, to facilitate free movement of goods inside the EU.
The previous core legal framework consisted of three directives:
Directive 90/385/EEC regarding active implantable medical devices
Directive 93/42/EEC regarding medical devices
Directive 98/79/EC regarding in vitro diagnostic medical devices (Until 2022, the In Vitro Diagnosis Regulation (IVDR) will replace the EU's current Directive on In-Vitro Diagnostic (98/79/EC)).
They aim at ensuring a high level of protection of human health and safety and the good functioning of the Single Market. These three main directives have been supplemented over time by several modifying and implementing directives, including the last technical revision brought about by Directive 2007/47 EC.
The government of each Member State must appoint a competent authority responsible for medical devices. The competent authority (CA) is a body with authority to act on behalf of the member state to ensure that member state government transposes requirements of medical device directives into national law and applies them. The CA reports to the minister of health in the member state. The CA in one Member State has no jurisdiction in any other member state, but exchanges information and tries to reach common positions.
In the UK, for example, the Medicines and Healthcare products Regulatory Agency (MHRA) acted as a CA. In Italy it is the Ministero Salute (Ministry of Health) Medical devices must not be mistaken with medicinal products. In the EU, all medical devices must be identified with the CE mark. The conformity of a medium or high risk medical device with relevant regulations is also assessed by an external entity, the Notified Body, before it can be placed on the market.
In September 2012, the European Commission proposed new legislation aimed at enhancing safety, traceability, and transparency. The regulation was adopted in 2017.
The currct core legal framework consists of two regulations, replacing the previous three directives:
The Medical Devices Regulation (MDR (EU) 2017/745)
The In Vitro Diagnostic medical devices regulation (IVDR (EU) 2017/746)
The two regulations are supplemented by several guidances developed by the Medical Devices Coordination Group (MDCG).
Japan
Article 2, Paragraph 4, of the Pharmaceutical Affairs Law (PAL) defines medical devices as "instruments and apparatus intended for use in diagnosis, cure or prevention of diseases in humans or other animals; intended to affect the structure or functions of the body of man or other animals."
Rest of the world
Canada
The term medical device, as defined in the Food and Drugs Act, is "any article, instrument, apparatus or contrivance, including any component, part or accessory thereof, manufactured, sold or represented for use in: the diagnosis, treatment, mitigation or prevention of a disease, disorder or abnormal physical state, or its symptoms, in a human being; the restoration, correction or modification of a body function or the body structure of a human being; the diagnosis of pregnancy in a human being; or the care of a human being during pregnancy and at and after the birth of a child, including the care of the child. It also includes a contraceptive device but does not include a drug."
The term covers a wide range of health or medical instruments used in the treatment, mitigation, diagnosis or prevention of a disease or abnormal physical condition. Health Canada reviews medical devices to assess their safety, effectiveness, and quality before authorizing their sale in Canada. According to the Act, medical device does not include any device that is intended for use in relation to animals.
India
India has introduced National Medical Device Policy 2023. However, certain medical devices are notified as DRUGS under the Drugs & Cosmetics Act. Section 3 (b) (iv) relating to definition of "drugs" holds that "Devices intended for internal or external use in the diagnosis, treatment, mitigation or prevention of disease or disorder in human beings or animals" are also drugs. As of April 2022, 14 classes of devices are classified as drugs.
Regulation and oversight
Risk classification
The regulatory authorities recognize different classes of medical devices based on their potential for harm if misused, design complexity, and their use characteristics. Each country or region defines these categories in different ways. The authorities also recognize that some devices are provided in combination with drugs, and regulation of these combination products takes this factor into consideration.
Classifying medical devices based on their risk is essential for maintaining patient and staff safety while simultaneously facilitating the marketing of medical products. By establishing different risk classifications, lower risk devices, for example, a stethoscope or tongue depressor, are not required to undergo the same level of testing that higher risk devices such as artificial pacemakers undergo. Establishing a hierarchy of risk classification allows regulatory bodies to provide flexibility when reviewing medical devices.
Classification by region
United States
Under the Food, Drug, and Cosmetic Act, the U.S. Food and Drug Administration recognizes three classes of medical devices, based on the level of control necessary to assure safety and effectiveness.
Class I
Class II
Class III
The classification procedures are described in the Code of Federal Regulations, Title 21, part 860 (usually known as 21 CFR 860).
Class I devices are subject to the least regulatory control and are not intended to help support or sustain life or be substantially important in preventing impairment to human health, and may not present an unreasonable risk of illness or injury. Examples of Class I devices include elastic bandages, examination gloves, and hand-held surgical instruments.
Class II devices are subject to special labeling requirements, mandatory performance standards and postmarket surveillance. Examples of Class II devices include acupuncture needles, powered wheelchairs, infusion pumps, air purifiers, surgical drapes, stereotaxic navigation systems, and surgical robots.
Class III devices are usually those that support or sustain human life, are of substantial importance in preventing impairment of human health, or present a potential, unreasonable risk of illness or injury and require premarket approval. Examples of Class III devices include implantable pacemakers, pulse generators, HIV diagnostic tests, automated external defibrillators, and endosseous implants.
European Union (EU) and European Free Trade Association (EFTA)
The classification of medical devices in the European Union is outlined in Article IX of the Council Directive 93/42/EEC and Annex VIII of the EU medical device regulation. There are basically four classes, ranging from low risk to high risk, Classes I, IIa, IIb, and III (this excludes in vitro diagnostics including software, which fall in four classes: from A (lowest risk) to D (highest risk)):
Class I Devices: Non-invasive, everyday devices or equipment. Class I devices are generally low risk and can include bandages, compression hosiery, or walking aids. Such devices require only for the manufacturer to complete a Technical File.
Class Is Devices: Class Is devices are similarly non-invasive devices, however this sub-group extends to include sterile devices. Examples of Class Is devices include stethoscopes, examination gloves, colostomy bags, or oxygen masks. These devices also require a technical file, with the added requirement of an application to a European Notified Body for certification of manufacturing in conjunction with sterility standards.
Class Im Devices: This refers chiefly to similarly low-risk measuring devices. Included in this category are: thermometers, droppers, and non-invasive blood pressure measuring devices. Once again the manufacturer must provide a technical file and be certified by a European Notified Body for manufacturing in accordance with metrology regulations.
Class IIa Devices: Class IIa devices generally constitute low to medium risk and pertain mainly to devices installed within the body in the short term. Class IIa devices are those which are installed within the body for only between 60 minutes and 30 days. Examples include hearing-aids, blood transfusion tubes, and catheters. Requirements include technical files and a conformity test carried out by a European Notified Body.
Class IIb Devices: Slightly more complex than IIa devices, class IIb devices are generally medium to high risk and will often be devices installed within the body for periods of 30 days or longer. Examples include ventilators and intensive care monitoring equipment. Identical compliance route to Class IIa devices with an added requirement of a device type examination by a Notified Body.
Class III Devices: Class III devices are strictly high risk devices. Examples include balloon catheters, prosthetic heart valves, pacemakers, etc. The steps to approval here include a full quality assurance system audit, along with examination of both the device's design and the device itself by a European Notified Body.
The authorization of medical devices is guaranteed by a Declaration of Conformity. This declaration is issued by the manufacturer itself, but for products in Class Is, Im, Ir, IIa, IIb or III, it must be verified by a Certificate of Conformity issued by a Notified Body. A Notified Body is a public or private organisation that has been accredited to validate the compliance of the device to the European Directive. Medical devices that pertain to class I (on condition they do not require sterilization or do not measure a function) can be marketed purely by self-certification.
The European classification depends on rules that involve the medical device's duration of body contact, invasive character, use of an energy source, effect on the central circulation or nervous system, diagnostic impact, or incorporation of a medicinal product. Certified medical devices should have the CE mark on the packaging, insert leaflets, etc.. These packagings should also show harmonised pictograms and EN standardised logos to indicate essential features such as instructions for use, expiry date, manufacturer, sterile, do not reuse, etc.
In November 2018, the Federal Administrative Court of Switzerland decided that the "Sympto" app, used to analyze a woman's menstrual cycle, was a medical device because it calculates a fertility window for each woman using personal data. The manufacturer, Sympto-Therm Foundation, argued that this was a didactic, not a medical process. the court laid down that an app is a medical device if it is to be used for any of the medical purposes provided by law, and creates or modifies health information by calculations or comparison, providing information about an individual patient.
Japan
Medical devices (excluding in vitro diagnostics) in Japan are classified into four classes based on risk:
Classes I and II distinguish between extremely low and low risk devices. Classes III and IV, moderate and high risk respectively, are highly and specially controlled medical devices. In vitro diagnostics have three risk classifications.
Rest of the world
For the remaining regions in the world, the risk classifications are generally similar to the United States, European Union, and Japan or are a variant combining two or more of the three countries' risk classifications.
ASEAN
The ASEAN Medical Device Directive (AMDD) has been adopted by several southeast Asian countries. The nations are at varying stages of adopting and implementing the Directive. The AMDD classification is risk-based and defines four levels: A - Low Risk, B - Low to Moderate Risk, C - Moderate – High Risk, and D - High Risk.
Australia
The classification of medical devices in Australia is outlined in section 41BD of the Therapeutic Goods Act 1989 and Regulation 3.2 of the Therapeutic Goods Regulations 2002, under control of the Therapeutic Goods Administration. Similarly to the EU classification, they rank in several categories, by order of increasing risk and associated required level of control. Various rules identify the device's category
Canada
The Medical Devices Bureau of Health Canada recognizes four classes of medical devices based on the level of control necessary to assure the safety and effectiveness of the device. Class I devices present the lowest potential risk and do not require a licence. Class II devices require the manufacturer's declaration of device safety and effectiveness, whereas Class III and IV devices present a greater potential risk and are subject to in-depth scrutiny. A guidance document for device classification is published by Health Canada.
Canadian classes of medical devices correspond to the European Council Directive 93/42/EEC (MDD) devices:
Class I (Canada) generally corresponds to Class I (ECD)
Class II (Canada) generally corresponds to Class IIa (ECD)
Class III (Canada) generally corresponds to Class IIb (ECD)
Class IV (Canada) generally corresponds to Class III (ECD)
Examples include surgical instruments (Class I), contact lenses and ultrasound scanners (Class II), orthopedic implants and hemodialysis machines (Class III), and cardiac pacemakers (Class IV).
India
Medical devices in India are regulated by Central Drugs Standard Control Organisation (CDSCO). Medical devices under the Medical Devices Rules, 2017 are classified as per Global Harmonization Task Force (GHTF) based on associated risks.
The CDSCO classifications of medical devices govern alongside the regulatory approval and registration by the CDSCO is under the DCGI. Every single medical device in India pursues a regulatory framework that depends on the drug guidelines under the Drug and Cosmetics Act (1940) and the Drugs and Cosmetics runs under 1945. CDSCO classification for medical devices has a set of risk classifications for numerous products planned for notification and guidelines as medical devices.
Iran
Iran produces about 2,000 types of medical devices and medical supplies, such as appliances, dental supplies, disposable sterile medical items, laboratory machines, various biomaterials and dental implants. 400 Medical products are produced at the C and D risk class with all of them licensed by the Iranian Health Ministry in terms of safety and performance based on EU-standards.
Some Iranian medical devices are produced according to the European Union standards.
Some producers in Iran export medical devices and supplies which adhere to European Union standards to applicant countries, including 40 Asian and European countries.
Some Iranian producers export their products to foreign countries.
United Kingdom
Following Brexit, the UK medical device regulation was closely aligned with the EU medical device regulation, including classification. The regulation 7 of the Medical Devices Regulations 2002 (SI 2002 No 618, as amended) (UK medical devices regulations), classified general medical devices into four classes of increasing levels of risk: Class I, IIa, IIb or III in accordance with criteria in the UK medical devices regulations, Annex IX (as modified by Schedule 2A to the UK medical devices regulations).
Validation and verification
Validation and verification of medical devices ensure that they fulfil their intended purpose. Validation or verification is generally needed when a health facility acquires a new device to perform medical tests.
The main difference between the two is that validation is focused on ensuring that the device meets the needs and requirements of its intended users and the intended use environment, whereas verification is focused on ensuring that the device meets its specified design requirements.
Standardization and regulatory concerns
The ISO standards for medical devices are covered by ICS 11.100.20 and 11.040.01. The quality and risk management regarding the topic for regulatory purposes is convened by ISO 13485 and ISO 14971. ISO 13485:2016 is applicable to all providers and manufacturers of medical devices, components, contract services and distributors of medical devices. The standard is the basis for regulatory compliance in local markets, and most export markets. Additionally, ISO 9001:2008 sets precedence because it signifies that a company engages in the creation of new products. It requires that the development of manufactured products have an approval process and a set of rigorous quality standards and development records before the product is distributed. Further standards are IEC 60601-1 which is for electrical devices (mains-powered as well as battery powered), EN 45502-1 which is for Active implantable medical devices, and IEC 62304 for medical software. The US FDA also published a series of guidances for industry regarding this topic against 21 CFR 820 Subchapter H—Medical Devices. Subpart B includes quality system requirements, an important component of which are design controls (21 CFR 820.30). To meet the demands of these industry regulation standards, a growing number of medical device distributors are putting the complaint management process at the forefront of their quality management practices. This approach further mitigates risks and increases visibility of quality issues.
Starting in the late 1980s, the FDA increased its involvement in reviewing the development of medical device software. The precipitant for change was a radiation therapy device (Therac-25) that overdosed patients because of software coding errors. FDA is now focused on regulatory oversight on medical device software development process and system-level testing.
A 2011 study by Dr. Diana Zuckerman and Paul Brown of the National Center for Health Research, and Dr. Steven Nissen of the Cleveland Clinic, published in the Archives of Internal Medicine, showed that most medical devices recalled in the last five years for "serious health problems or death" had been previously approved by the FDA using the less stringent, and cheaper, 510(k) process. In a few cases, the devices had been deemed so low-risk that they did not they did not undergo any FDA regulatory review. Of the 113 devices recalled, 35 were for cardiovascular issues. This study was the topic of Congressional hearings re-evaluating FDA procedures and oversight.
A 2014 study by Dr. Diana Zuckerman, Paul Brown, and Dr. Aditi Das of the National Center for Health Research, published in JAMA Internal Medicine, examined the scientific evidence that is publicly available about medical implants that were cleared by the FDA 510(k) process from 2008 to 2012. They found that scientific evidence supporting "substantial equivalence" to other devices already on the market was required by law to be publicly available, but the information was available for only 16% of the randomly selected implants, and only 10% provided clinical data. Of the more than 1,100 predicate implants that the new implants were substantially equivalent to, only 3% had any publicly available scientific evidence, and only 1% had clinical evidence of safety or effectiveness. The researchers concluded that publicly available scientific evidence on implants was needed to protect the public health.
In 2014–2015, a new international agreement, the Medical Device Single Audit Program (MDSAP), was put in place with five participant countries: Australia, Brazil, Canada, Japan, and the United States. The aim of this program was to "develop a process that allows a single audit, or inspection to ensure the medical device regulatory requirements for all five countries are satisfied".
In 2017, a study by Dr. Jay Ronquillo and Dr. Diana Zuckerman published in the peer-reviewed policy journal Milbank Quarterly found that electronic health records and other device software were recalled due to life-threatening flaws. The article pointed out the lack of safeguards against hacking and other cybersecurity threats, stating "current regulations are necessary but not sufficient for ensuring patient safety by identifying and eliminating dangerous defects in software currently on the market". They added that legislative changes resulting from the law entitled the 21st Century Cures Act "will further deregulate health IT, reducing safeguards that facilitate the reporting and timely recall of flawed medical software that could harm patients".
A study by Dr. Stephanie Fox-Rawlings and colleagues at the National Center for Health Research, published in 2018 in the policy journal Milbank Quarterly, investigated whether studies reviewed by the FDA for high-risk medical devices are proven safe and effective for women, minorities, or patients over 65 years of age. The law encourages patient diversity in clinical trials submitted to the FDA for review, but does not require it. The study determined that most high-risk medical devices are not tested and analyzed to ensure that they are safe and effective for all major demographic groups, particularly racial and ethnic minorities and people over 65. Therefore, they do not provide information about safety or effectiveness that would help patients and physicians make well informed decisions.
In 2018, an investigation involving journalists across 36 countries coordinated by the International Consortium of Investigative Journalists (ICIJ) prompted calls for reform in the United States, particularly around the 510(k) substantial equivalence process; the investigation prompted similar calls in the UK and Europe Union.
Packaging standards
Medical device packaging is highly regulated. Often medical devices and products are sterilized in the package.
Sterility must be maintained throughout distribution to allow immediate use by physicians. A series of special packaging tests measure the ability of the package to maintain sterility. Relevant standards include:
ASTM F2097 – Standard Guide for Design and Evaluation of Primary Flexible Packaging for Medical Products
ASTM F2475-11 – Standard Guide for Biocompatibility Evaluation of Medical Device Packaging Materials
EN 868 Packaging materials and systems for medical devices to be sterilized, General requirements and test methods
ISO 11607 Packaging for terminally sterilized medical devices
Package testing is part of a quality management system including verification and validation. It is important to document and ensure that packages meet regulations and end-use requirements. Manufacturing processes must be controlled and validated to ensure consistent performance. EN ISO 15223-1 defines symbols that can be used to convey important information on packaging and labeling.
Biocompatibility standards
ISO 10993 - Biological Evaluation of Medical Devices
Cleanliness standards
Medical device cleanliness has come under greater scrutiny since 2000, when Sulzer Orthopedics recalled several thousand metal hip implants that contained a manufacturing residue. Based on this event, ASTM established a new task group (F04.15.17) for established test methods, guidance documents, and other standards to address cleanliness of medical devices. This task group has issued two standards for permanent implants to date: 1. ASTM F2459: Standard test method for extracting residue from metallic medical components and quantifying via gravimetric analysis 2. ASTM F2847: Standard Practice for Reporting and Assessment of Residues on Single Use Implants 3. ASTM F3172: Standard Guide for Validating Cleaning Processes Used During the Manufacture of Medical Devices
In addition, the cleanliness of re-usable devices has led to a series of standards, including:
ASTM E2314: Standard Test Method for Determination of Effectiveness of Cleaning Processes for Reusable Medical Instruments Using a Microbiologic Method (Simulated Use Test)"
ASTM D7225: Standard Guide for Blood Cleaning Efficiency of Detergents and Washer-Disinfectors
ASTM F3208: Standard Guide for Selecting Test Soils for Validation of Cleaning Methods for Reusable Medical Devices
The ASTM F04.15.17 task group is working on several new standards that involve designing implants for cleaning, selection and testing of brushes for cleaning reusable devices, and cleaning assessment of medical devices made by additive manufacturing. Additionally, the FDA is establishing new guidelines for reprocessing reusable medical devices, such as orthoscopic shavers, endoscopes, and suction tubes. New research was published in ACS Applied Interfaces and Material to keep Medical Tools pathogen free.
Safety standards
Design, prototyping, and product development
Medical device manufacturing requires a level of process control according to the classification of the device. Higher risk; more controls. When in the initial R&D phase, manufacturers are now beginning to design for manufacturability. This means products can be more precision-engineered to for production to result in shorter lead times, tighter tolerances and more advanced specifications and prototypes. These days, with the aid of CAD or modelling platforms, the work is now much faster, and this can act also as a tool for strategic design generation as well as a marketing tool.
Failure to meet cost targets will lead to substantial losses for an organisation. In addition, with global competition, the R&D of new devices is not just a necessity, it is an imperative for medical device manufacturers. The realisation of a new design can be very costly, especially with the shorter product life cycle. As technology advances, there is typically a level of quality, safety and reliability that increases exponentially with time.
For example, initial models of the artificial cardiac pacemaker were external support devices that transmits pulses of electricity to the heart muscles via electrode leads on the chest. The electrodes contact the heart directly through the chest, allowing stimulation pulses to pass through the body. Recipients of this typically developed an infection at the entrance of the electrodes, which led to the subsequent trial of the first internal pacemaker, with electrodes attached to the myocardium by thoracotomy. Future developments led to the isotope-power source that would last for the lifespan of the patient.
Software
Mobile medical applications
With the rise of smartphone usage in the medical space, in 2013, the FDA issued to regulate mobile medical applications and protect users from their unintended use, soon followed by European and other regulatory agencies. This guidance distinguishes the apps subjected to regulation based on the marketing claims of the apps. Incorporation of the guidelines during the development phase of such apps can be considered as developing a medical device; the regulations have to adapt and propositions for expedite approval may be required due to the nature of 'versions' of mobile application development.
On September 25, 2013, the FDA released a draft guidance document for regulation of mobile medical applications, to clarify what kind of mobile apps related to health would not be regulated, and which would be.
Cybersecurity
Medical devices such as pacemakers, insulin pumps, operating room monitors, defibrillators, and surgical instruments, including deep-brain stimulators, can incorporate the ability to transmit vital health information from a patient's body to medical professionals. Some of these devices can be remotely controlled. This has engendered concern about privacy and security issues, human error, and technical glitches with this technology. While only a few studies have looked at the susceptibility of medical devices to hacking, there is a risk. In 2008, computer scientists proved that pacemakers and defibrillators can be hacked wirelessly via radio hardware, an antenna, and a personal computer. These researchers showed they could shut down a combination heart defibrillator and pacemaker and reprogram it to deliver potentially lethal shocks or run out its battery. Jay Radcliff, a security researcher interested in the security of medical devices, raised fears about the safety of these devices. He shared his concerns at the Black Hat security conference. Radcliff fears that the devices are vulnerable and has found that a lethal attack is possible against those with insulin pumps and glucose monitors. Some medical device makers downplay the threat from such attacks and argue that the demonstrated attacks have been performed by skilled security researchers and are unlikely to occur in the real world. At the same time, other makers have asked software security experts to investigate the safety of their devices. As recently as June 2011, security experts showed that by using readily available hardware and a user manual, a scientist could both tap into the information on the system of a wireless insulin pump in combination with a glucose monitor. With the PIN of the device, the scientist could wirelessly control the dosage of the insulin. Anand Raghunathan, a researcher in this study, explains that medical devices are getting smaller and lighter so that they can be easily worn. The downside is that additional security features would put an extra strain on the battery and size and drive up prices. Dr. William Maisel offered some thoughts on the motivation to engage in this activity. Motivation to do this hacking might include acquisition of private information for financial gain or competitive advantage; damage to a device manufacturer's reputation; sabotage; intent to inflict financial or personal injury or just satisfaction for the attacker. Researchers suggest a few safeguards. One would be to use rolling codes. Another solution is to use a technology called "body-coupled communication" that uses the human skin as a wave guide for wireless communication. On 28 December 2016, the US Food and Drug Administration released its recommendations that are not legally enforceable for how medical device manufacturers should maintain the security of Internet-connected devices.
Similar to hazards, cybersecurity threats and vulnerabilities cannot be eliminated but must be managed and reduced to a reasonable level. When designing medical devices, the tier of cybersecurity risk should be determined early in the process in order to establish a cybersecurity vulnerability and management approach (including a set of cybersecurity design controls). The medical device design approach employed should be consistent with the NIST Cybersecurity Framework for managing cybersecurity-related risks.
In August 2013, the FDA released over 20 regulations aiming to improve the security of data in medical devices, in response to the growing risks of limited cybersecurity.
Artificial intelligence
The number of approved medical devices using artificial intelligence or machine learning (AI/ML) is increasing. As of 2020, there were several hundred AI/ML medical devices approved by the US FDA or CE-marked devices in Europe. Most AI/ML devices focus upon radiology. As of 2020, there was no specific regulatory pathway for AI/ML-based medical devices in the US or Europe. However, in January 2021, the FDA published a proposed regulatory framework for AI/ML-based software, and the EU medical device regulation which replaces the EU Medical Device Directive in May 2021, defines regulatory requirements for medical devices, including AI/ML software.
Medical equipment
Medical equipment (also known as armamentarium) is designed to aid in the diagnosis, monitoring or treatment of medical conditions.
Types
There are several basic types:
Diagnostic equipment includes medical imaging machines, used to aid in diagnosis. Examples are ultrasound and MRI machines, PET and CT scanners, and x-ray machines.
Treatment equipment includes infusion pumps, medical lasers and LASIK surgical machines.
Life support equipment is used to maintain a patient's bodily function. This includes medical ventilators, incubators, anaesthetic machines, heart-lung machines, ECMO, and dialysis machines.
Medical monitors allow medical staff to measure a patient's medical state. Monitors may measure patient vital signs and other parameters including ECG, EEG, and blood pressure.
Medical laboratory equipment automates or helps analyze blood, urine, genes, and dissolved gases in the blood.
Diagnostic medical equipment may also be used in the home for certain purposes, e.g. for the control of diabetes mellitus, such as in the case of continuous glucose monitoring.
Therapeutic: physical therapy machines like continuous passive range of motion (CPM) machines
Air purifying equipment may be used in the periphery of the operating room or at point sources including near the surgical site for the removal of surgical plume.
The identification of medical devices has been recently improved by the introduction of Unique Device Identification (UDI) and standardised naming using the Global Medical Device Nomenclature (GMDN) which have been endorsed by the International Medical Device Regulatory Forum (IMDRF).
A biomedical equipment technician (BMET) is a vital component of the healthcare delivery system. Employed primarily by hospitals, BMETs are the people responsible for maintaining a facility's medical equipment. BMET mainly act as an interface between doctor and equipment.
Medical equipment donation
There are challenges surrounding the availability of medical equipment from a global health perspective, with low-resource countries unable to obtain or afford essential and life-saving equipment. In these settings, well-intentioned equipment donation from high- to low-resource settings is a frequently used strategy to address this through individuals, organisations, manufacturers and charities. However, issues with maintenance, availability of biomedical equipment technicians (BMET), supply chains, user education and the appropriateness of donations means these frequently fail to deliver the intended benefits. The WHO estimates that 95% of medical equipment in low- and middle-income countries (LMICs) is imported and 80% of it is funded by international donors or foreign governments. While up to 70% of medical equipment in sub-Saharan Africa is donated, only 10%–30% of donated equipment becomes operational. A review of current practice and guidelines for the donation of medical equipment for surgical and anaesthesia care in LMICs has demonstrated a high level of complexity within the donation process and numerous shortcomings. Greater collaboration and planning between donors and recipients is required together with evaluation of donation programs and concerted advocacy to educate donors and recipients on existing equipment donation guidelines and policies.
The circulation of medical equipment is not limited to donations. The rise of reuse and recycle-based solutions, where gently-used medical equipment is donated and redistributed to communities in need, is another form of equipment distribution. An interest in reusing and recycling emerged in the 1980s when the potential health hazards of medical waste on the East Coast beaches became highlighted by the media. Connecting the large demand for medical equipment and single-use medical devices, with a need for waste reduction, as well as the problem of unequal access for low-income communities led to the Congress enacting the Medical Waste Tracking Act of 1988. Medical equipment can be donated either by governments or non-governmental organizations, domestic or international. Donated equipment ranges from bedside assistance to radiological equipment.
Medical equipment donation has come under scrutiny with regard to donated-device failure and loss of warranty in the case of previous-ownership. Most medical devices and production company warranties to do not extend to reused or donated devices, or to devices donated by initial owners/patients. Such reuse raises matters of patient autonomy, medical ethics, and legality. Such concerns conflict with the importance of equal access to healthcare resources, and the goal of serving the greatest good for the greatest number.
Academic resources
Medical & Biological Engineering & Computing journal
Expert Review of Medical Devices journal
University-based research packaging institutes
University of Minnesota - Medical Devices Center (MDC)
University of Strathclyde - Strathclyde Institute of Medical Devices (SIMD)
Flinders University - Medical Device Research Institute (MDRI)
Michigan State University - School of Packaging (SoP)
IIT Bombay - Biomedical Engineering and Technology (incubation) Centre (BETiC)
| Biology and health sciences | General concepts | Health |
1363657 | https://en.wikipedia.org/wiki/Rad%20%28radiation%20unit%29 | Rad (radiation unit) | The rad is a unit of absorbed radiation dose, defined as 1 rad = 0.01 Gy = 0.01 J/kg. It was originally defined in CGS units in 1953 as the dose causing 100 ergs of energy to be absorbed by one gram of matter. The material absorbing the radiation can be human tissue, air, water, or any other substance.
It has been replaced by the gray (symbol Gy) in SI derived units, but is still used in the United States, although this is "strongly discouraged" in Chapter 5.2 of the Guide to the SI, which was written and published by the U.S. National Institute of Standards and Technology. However, the numerically equivalent SI unit submultiple, the centigray (symbol cGy), is widely used to report absorbed doses within radiotherapy. The roentgen, used to quantify the radiation exposure, may be related to the corresponding absorbed dose by use of the F-factor.
Health effects
A dose of under 100 rad will typically produce no immediate symptoms other than blood changes. A dose of 100 to 200 rad delivered to the entire body in less than a day may cause acute radiation syndrome (ARS), but is usually not fatal. Doses of 200 to 1,000 rad delivered in a few hours will cause serious illness, with poor prognosis at the upper end of the range. Whole body doses of more than 1,000 rad are almost invariably fatal. Therapeutic doses of radiation therapy are often given and tolerated well even at higher doses to treat discrete, well-defined anatomical structures. The same dose given over a longer period of time is less likely to cause ARS. Dose thresholds are about 50% higher for dose rates of 20 rad/h, and even higher for lower dose rates.
The International Commission on Radiological Protection maintains a model of health risks as a function of absorbed dose and other factors. That model calculates an effective radiation dose, measured in units of rem, which is more representative of the stochastic risk than the absorbed dose in rad. In most power plant scenarios, where the radiation environment is dominated by X- or gamma rays applied uniformly to the whole body, 1 rad of absorbed dose gives 1 rem of effective dose. In other situations, the effective dose in rem might be thirty times higher or thousands of times lower than the absorbed dose in rad.
History
In the 1930s the roentgen was the most commonly used unit of radiation exposure. This unit is obsolete and no longer clearly defined. One roentgen deposits 0.877 rad in dry air, 0.96 rad in soft tissue, or anywhere from 1 to more than 4 rad in bone depending on the beam energy. These conversions to absorbed energy all depend on the ionizing energy of a standard medium, which is ambiguous in the latest NIST definition. Even where the standard medium is fully defined, the ionizing energy is often not precisely known.
In 1940, British physicist Louis Harold Gray, who had been studying the effect of neutron damage on human tissue, together with William Valentine Mayneord and John Read published a paper in which a unit of measure, dubbed the "gram roentgen" (symbol: gr) defined as "that amount of neutron radiation which produces an increment in energy in unit volume of tissue equal to the increment of energy produced in unit volume of water by one roentgen of radiation" was proposed. This unit was found to be equivalent to 88 ergs in air. It marked a shift towards measurements based on energy rather than charge.
The Röntgen equivalent physical (rep), introduced by Herbert Parker in 1945, was the absorbed energetic dose to tissue before factoring in relative biological effectiveness. The rep has variously been defined as 83 or 93 ergs per gram of tissue (8.3/9.3 mGy) or per cc of tissue.
In 1953 the ICRU recommended the rad, equal to 100 erg/g as a new unit of absorbed radiation, but then promoted a switch to the gray in the 1970s.
The International Committee for Weights and Measures (CIPM) has not accepted the use of the rad. From 1977 to 1998, the US NIST's translations of the SI brochure stated that the CIPM had temporarily accepted the use of the rad (and other radiology units) with SI units since 1969. However, the only related CIPM decisions shown in the appendix are with regards to the curie in 1964 and the radian (symbol: rad) in 1960. The NIST brochures redefined the rad as 0.01 Gy. The CIPM's current SI brochure excludes the rad from the tables of non-SI units accepted for use with the SI. The US NIST clarified in 1998 that it was providing its own interpretations of the SI system, whereby it accepted the rad for use in the US with the SI, while recognizing that the CIPM did not. NIST recommends defining the rad in relation to SI units in every document where this unit is used. Nevertheless, use of the rad remains widespread in the US, where it is still an industry standard. Although the United States Nuclear Regulatory Commission still permits the use of the units curie, rad, and rem alongside SI units, the European Union required that its use for "public health ... purposes" be phased out by 31 December 1985.
Radiation-related quantities
The following table shows radiation quantities in SI and non-SI units:
| Physical sciences | Absorbed dose | Basics and measurement |
1364340 | https://en.wikipedia.org/wiki/Cape%20%28geography%29 | Cape (geography) | In geography, a cape is a headland, peninsula or promontory extending into a body of water, usually a sea. A cape usually represents a marked change in trend of the coastline, often making them important landmarks in sea navigation. This also makes them prone to natural forms of erosion, mainly tidal actions, resulting in a relatively short geological lifespan.
Formation
Capes can be formed by glaciers, volcanoes, and changes in sea level. Erosion plays a large role in each of these methods of formation. Coastal erosion by waves and currents can create capes by wearing away softer rock and leaving behind harder rock formations. Movements of the Earth's crust can uplift land, forming capes. For example, the Cape of Good Hope was formed by tectonic forces. Volcanic eruptions can create capes by depositing lava that solidifies into new landforms. Cape Verde, (also known as Cabo Verde) is an example of a volcanic cape. Glaciers can carve out capes by eroding the landscape as they advance and retreat. Cape Cod in the United States was formed by glacial activity during the last Ice Age.
Importance in navigation
Capes (and other headlands) are conspicuous visual landmarks along a coast, and sailors have relied on them for navigation since antiquity. The Greeks and Romans considered some to be sacred capes and erected temples to the sea god nearby.
Greek peripli describe capes and other headlands a sailor will encounter along a route. The Periplus of Pseudo-Scylax, for instance, illustrates a clockwise journey around Sicily using three capes that define its triangular shape: Cape Peloro in the northeast, Cape Pachynus in the southeast, and Cape Lilybaeum in the west. Sicily itself was referred to as Trinacria (or Three Capes) in antiquity.
Homer's works reference a number of capes to describe journeys around the Mediterranean Sea. Menelaus, Agamemnon, and Odysseus each faced peril at the notoriously dangerous Cape Malea at the southeastern tip of the Peloponnese. Menelaus navigated via Cape Sounion on his way home from Troy, and Nestor stopped at Cape Geraestus (now Cape Mandelo) on Euboea to give offerings at the altar to Poseidon there.
Cape Gelidonya (then known as Chelidonia) on the coast of Turkey served as a bearing aid for ships heading to the Egyptian port of Canopus, directly to the south. Cape Sidero on the eastern tip of Crete was a waypoint for Jason and the Argonauts returning from Libya as well as for Paul the Apostle as he traveled from Caesarea to Rome.
The three great capes (Africa's Cape of Good Hope, Australia's Cape Leeuwin, and South America's Cape Horn) defined the traditional clipper route between Europe and the Far East, Australia and New Zealand. They continue to be important landmarks in ocean yacht racing.
List of capes
Antarctica
Cape Ann
Cape May (Antarctica), McMurdo Sound
Chile
Cape Horn, Chile
India
Cape Comorin, India
United States
Cape Ann, Massachusetts
Cape Cod, Massachusetts
Cape May, New Jersey
Cape Charles, Virginia
Cape Henry, Virginia
Cape Hatteras, North Carolina
Cape Lookout, North Carolina
Cape Fear, North Carolina
Cape Canaveral, Florida
Cape Canaveral Space Force Station, a launch station of the US Space Force
Cape Coral, Florida
Cape Rosier, Maine
South Africa
Cape of Good Hope, a headland on the southwest coast of South Africa, when referred to as the Cape, a metonym for:
Dutch Cape Colony, a colony of the Dutch East India company
Cape Colony, a British colony in South Africa that replaced the Dutch Cape Colony
Cape Province, a former province of South Africa formed from the Cape Colony
Cape Town, a city in South Africa, and surrounding areas
Gallery
| Physical sciences | Oceanic and coastal landforms | Earth science |
18949797 | https://en.wikipedia.org/wiki/File%20sharing | File sharing | File sharing is the practice of distributing or providing access to digital media, such as computer programs, multimedia (audio, images and video), documents or electronic books. Common methods of storage, transmission and dispersion include removable media, centralized servers on computer networks, Internet-based hyperlinked documents, and the use of distributed peer-to-peer networking.
File sharing technologies, such as BitTorrent, are integral to modern media piracy, as well as the sharing of scientific data and other free content.
History
Files were first exchanged on removable media. Computers were able to access remote files using filesystem mounting, bulletin board systems (1978), Usenet (1979), and FTP servers (1970's). Internet Relay Chat (1988) and Hotline (1997) enabled users to communicate remotely through chat and to exchange files. The mp3 encoding, which was standardized in 1991 and substantially reduced the size of audio files, grew to widespread use in the late 1990s. In 1998, MP3.com and Audiogalaxy were established, the Digital Millennium Copyright Act was unanimously passed, and the first mp3 player devices were launched.
In June 1999, Napster was released as an unstructured centralized peer-to-peer system, requiring a central server for indexing and peer discovery. It is generally credited as being the first peer-to-peer file sharing system. In December 1999, Napster was sued by several recording companies and lost in A&M Records, Inc. v. Napster, Inc.. In the case of Napster, it has been ruled that an online service provider could not use the "transitory network transmission" safe harbor in the DMCA if they had control of the network with a server.
Gnutella, eDonkey2000, and Freenet were released in 2000, as MP3.com and Napster were facing litigation. Gnutella, released in March, was the first decentralized file-sharing network. In the Gnutella network, all connecting software was considered equal, and therefore the network had no central point of failure. In July, Freenet was released and became the first anonymity network. In September the eDonkey2000 client and server software was released.
In March 2001, Kazaa was released. Its FastTrack network was distributed, though, unlike Gnutella, it assigned more traffic to 'supernodes' to increase routing efficiency. The network was proprietary and encrypted, and the Kazaa team made substantial efforts to keep other clients such as Morpheus off of the FastTrack network. In October 2001, the MPAA and the RIAA filed a lawsuit against the developers of Kazaa, Morpheus and Grokster that would lead to the US Supreme Court's MGM Studios, Inc. v. Grokster, Ltd. decision in 2005.
Shortly after its loss in court, Napster was shut down to comply with a court order. This drove users to other P2P applications and file sharing continued its growth. The Audiogalaxy Satellite client grew in popularity, and the LimeWire client and BitTorrent protocol were released. Until its decline in 2004, Kazaa was the most popular file-sharing program despite bundled malware and legal battles in the Netherlands, Australia, and the United States. In 2002, a Tokyo district court ruling shut down File Rogue, and the Recording Industry Association of America (RIAA) filed a lawsuit that effectively shut down Audiogalaxy.
From 2002 through 2003, a number of BitTorrent services were established, including Suprnova.org, isoHunt, TorrentSpy, and The Pirate Bay. In September 2003, the RIAA began filing lawsuits against users of P2P file sharing networks such as Kazaa. As a result of such lawsuits, many universities added file sharing regulations in their school administrative codes (though some students managed to circumvent them during after school hours). Also in 2003, the MPAA started to take action against BitTorrent sites, leading to the shutdown of Torrentse and Sharelive in July 2003. With the shutdown of eDonkey in 2005, eMule became the dominant client of the eDonkey network. In 2006, police raids took down the Razorback2 eDonkey server and temporarily took down The Pirate Bay.
"The File Sharing Act was launched by Chairman Towns in 2009, this act prohibited the use of applications that allowed individuals to share federal information amongst one another. On the other hand, only specific file sharing applications were made available to federal computers" (the United States.Congress.House). In 2009, the Pirate Bay trial ended in a guilty verdict for the primary founders of the tracker. The decision was appealed, leading to a second guilty verdict in November 2010. In October 2010, Limewire was forced to shut down following a court order in Arista Records LLC v. Lime Group LLC but the Gnutella network remains active through open source clients like FrostWire and gtk-gnutella. Furthermore, multi-protocol file-sharing software such as MLDonkey and Shareaza adapted to support all the major file-sharing protocols, so users no longer had to install and configure multiple file-sharing programs.
On January 19, 2012, the United States Department of Justice shut down the popular domain of Megaupload (established 2005). The file sharing site has claimed to have over 50,000,000 people a day. Kim Dotcom (formerly Kim Schmitz) was arrested with three associates in New Zealand on January 20, 2012 and is awaiting extradition. The case involving the downfall of the world's largest and most popular file sharing site was not well received, with hacker group Anonymous bringing down several sites associated with the take-down. In the following days, other file sharing sites began to cease services; FileSonic blocked public downloads on January 22, with Fileserve following suit on January 23.
In 2021 a European Citizens' Initiative "Freedom to Share" started collecting signatures in order to get the European Commission to discuss (and eventually make rules) on this subject, which is controversial.
Techniques used for video sharing
From the early 2000s until the mid 2010s, online video streaming was usually based on the Adobe Flash Player. After more and more vulnerabilities in Adobe's flash became known, YouTube switched to HTML5 based video playback in January 2015.
Types
Peer-to-peer file sharing
Peer-to-peer file sharing is based on the peer-to-peer (P2P) application architecture. Shared files on the computers of other users are indexed on directory servers. P2P technology was used by popular services like Napster and LimeWire. The most popular protocol for P2P sharing is BitTorrent.
File sync and sharing services
Cloud-based file syncing and sharing services implement automated file transfers by updating files from a dedicated sharing directory on each user's networked devices. Files placed in this folder also are typically accessible through a website and mobile app and can be easily shared with other users for viewing or collaboration. Such services have become popular via consumer-oriented file hosting services such as Dropbox and Google Drive. With the rising need of sharing big files online easily, new open access sharing platforms have appeared, adding even more services to their core business (cloud storage, multi-device synchronization, online collaboration), such as ShareFile, Tresorit, WeTransfer, or Hightail.
rsync is a more traditional program released in 1996 which synchronizes files on a direct machine-to-machine basis.
Data synchronization in general can use other approaches to share files, such as distributed file systems, version control, or mirrors.
Academic file sharing
In addition to file sharing for the purposes of entertainment, academic file sharing has become a topic of increasing concern, as it is deemed to be a violation of academic integrity at many schools. Academic file sharing by companies such as Chegg and Course Hero has become a point of particular controversy in recent years. This has led some institutions to provide explicit guidance to students and faculty regarding academic integrity expectations relating to academic file sharing.
Public opinion of file sharing
In 2004, there were an estimated 70 million people participating in online file sharing. According to a CBS News poll in 2009, 58% of Americans who follow the file-sharing issue, considered it acceptable "if a person owns the music CD and shares it with a limited number of friends and acquaintances"; with 18- to 29-year-olds, this percentage reached as much as 70%.
In his survey of file-sharing culture, Caraway (2012) noted that 74.4% of participants believed musicians should accept file sharing as a means for promotion and distribution. This file-sharing culture was termed as cyber socialism, whose legalisation was not the expected cyber-utopia..
Economic impact
According to David Glenn, writing in The Chronicle of Higher Education, "A majority of economic studies have concluded that file-sharing hurts sales". A literature review by Professor Peter Tschmuck found 22 independent studies on the effects of music file sharing. "Of these 22 studies, 14 – roughly two-thirds – conclude that unauthorized downloads have a 'negative or even highly negative impact' on recorded music sales. Three of the studies found no significant impact while the remaining five found a positive impact."
A study by economists Felix Oberholzer-Gee and Koleman Strumpf in 2004 concluded that music file sharing's effect on sales was "statistically indistinguishable from zero". This research was disputed by other economists, most notably Stan Liebowitz, who said Oberholzer-Gee and Strumpf had made multiple assumptions about the music industry "that are just not correct." In June 2010, Billboard reported that Oberholzer-Gee and Strumpf had "changed their minds", now finding "no more than 20% of the recent decline in sales is due to sharing". However, citing Nielsen SoundScan as their source, the co-authors maintained that illegal downloading had not deterred people from being original. "In many creative industries, monetary incentives play a reduced role in motivating authors to remain creative. Data on the supply of new works are consistent with the argument that file-sharing did not discourage authors and publishers. Since the advent of file sharing, the production of music, books, and movies has increased sharply." Glenn Peoples of Billboard disputed the underlying data, saying "SoundScan's number for new releases in any given year represents new commercial titles, not necessarily new creative works." The RIAA likewise responded that "new releases" and "new creative works" are two separate things. "[T]his figure includes re-releases, new compilations of existing songs, and new digital-only versions of catalog albums. SoundScan has also steadily increased the number of retailers (especially non-traditional retailers) in their sample over the years, better capturing the number of new releases brought to market. What Oberholzer and Strumpf found was better ability to track new album releases, not greater incentive to create them."
A 2006 study prepared by Birgitte Andersen and Marion Frenz, published by Industry Canada, was "unable to discover any direct relationship between P2P file-sharing and CD purchases in Canada". The results of this survey were similarly criticized by academics and a subsequent revaluation of the same data by George R. Barker of the Australian National University reached the opposite conclusion. "In total, 75% of P2P downloaders responded that if P2P were not available they would have purchased either through paid sites only (9%), CDs only (17%) or through CDs and pay sites (49%). Only 25% of people say they would not have bought the music if it were not available on P2P for free." Barker thus concludes; "This clearly suggests P2P network availability is reducing music demand of 75% of music downloaders which is quite contrary to Andersen and Frenz's much published claim."
According to the 2017 paper "Estimating displacement rates of copyrighted content in the EU" by the European Commission, illegal usage increases game sales, stating "The overall conclusion is that for games, illegal online transactions induce more legal transactions."
Market dominance
A paper in the journal Management Science found that file-sharing decreased the chance of survival for low ranked albums on music charts and increased exposure to albums that were ranked high on the music charts, allowing popular and well-known artists to remain on the music charts more often. This hurt new and less-known artists while promoting the work of already popular artists and celebrities.
A more recent study that examined pre-release file-sharing of music albums, using BitTorrent software, also discovered positive impacts for "established and popular artists but not newer and smaller artists." According to Robert G. Hammond of North Carolina State University, an album that leaked one month early would see a modest increase in sales. "This increase in sales is small relative to other factors that have been found to affect album sales."
"File-sharing proponents commonly argue that file-sharing democratizes music consumption by 'levelling the playing field' for new/small artists relative to established/popular artists, by allowing artists to have their work heard by a wider audience, lessening the advantage held by established/popular artists in terms of promotional and other support. My results suggest that the opposite is happening, which is consistent with evidence on file-sharing behaviour."
Billboard cautioned that this research looked only at the pre-release period and not continuous file sharing following a release date. "The problem in believing piracy helps sales is deciding where to draw the line between legal and illegal ... Implicit in the study is the fact that both buyers and sellers are required in order for pre-release file sharing to have a positive impact on album sales. Without iTunes, Amazon, and Best Buy, file-sharers would be just file sharers rather than purchasers. If you carry out the 'file-sharing should be legal' argument to its logical conclusion, today's retailers will be tomorrow's file-sharing services that integrate with their respective cloud storage services."
Availability
Many argue that file-sharing has forced the owners of entertainment content to make it more widely available legally through fees or advertising on-demand on the internet. In a 2011 report by Sandvine showed that Netflix traffic had come to surpass that of BitTorrent.
Copyright issues
File sharing raises copyright issues and has led to many lawsuits. In the United States, some of these lawsuits have even reached the Supreme Court. For example, in MGM v. Grokster, the Supreme Court ruled that the creators of P2P networks can be held liable if their software is marketed as a tool for copyright infringement.
On the other hand, not all file sharing is illegal. Content in the public domain can be freely shared. Even works covered by copyright can be shared under certain circumstances. For example, some artists, publishers, and record labels grant the public a license for unlimited distribution of certain works, sometimes with conditions, and they advocate free content and file sharing as a promotional tool.
| Technology | Basics_4 | null |
18949896 | https://en.wikipedia.org/wiki/Computer%20cluster | Computer cluster | A computer cluster is a set of computers that work together so that they can be viewed as a single system. Unlike grid computers, computer clusters have each node set to perform the same task, controlled and scheduled by software. The newest manifestation of cluster computing is cloud computing.
The components of a cluster are usually connected to each other through fast local area networks, with each node (computer used as a server) running its own instance of an operating system. In most circumstances, all of the nodes use the same hardware and the same operating system, although in some setups (e.g. using Open Source Cluster Application Resources (OSCAR)), different operating systems can be used on each computer, or different hardware.
Clusters are usually deployed to improve performance and availability over that of a single computer, while typically being much more cost-effective than single computers of comparable speed or availability.
Computer clusters emerged as a result of the convergence of a number of computing trends including the availability of low-cost microprocessors, high-speed networks, and software for high-performance distributed computing. They have a wide range of applicability and deployment, ranging from small business clusters with a handful of nodes to some of the fastest supercomputers in the world such as IBM's Sequoia. Prior to the advent of clusters, single-unit fault tolerant mainframes with modular redundancy were employed; but the lower upfront cost of clusters, and increased speed of network fabric has favoured the adoption of clusters. In contrast to high-reliability mainframes, clusters are cheaper to scale out, but also have increased complexity in error handling, as in clusters error modes are not opaque to running programs.
Basic concepts
The desire to get more computing power and better reliability by orchestrating a number of low-cost commercial off-the-shelf computers has given rise to a variety of architectures and configurations.
The computer clustering approach usually (but not always) connects a number of readily available computing nodes (e.g. personal computers used as servers) via a fast local area network. The activities of the computing nodes are orchestrated by "clustering middleware", a software layer that sits atop the nodes and allows the users to treat the cluster as by and large one cohesive computing unit, e.g. via a single system image concept.
Computer clustering relies on a centralized management approach which makes the nodes available as orchestrated shared servers. It is distinct from other approaches such as peer-to-peer or grid computing which also use many nodes, but with a far more distributed nature.
A computer cluster may be a simple two-node system which just connects two personal computers, or may be a very fast supercomputer. A basic approach to building a cluster is that of a Beowulf cluster which may be built with a few personal computers to produce a cost-effective alternative to traditional high-performance computing. An early project that showed the viability of the concept was the 133-node Stone Soupercomputer. The developers used Linux, the Parallel Virtual Machine toolkit and the Message Passing Interface library to achieve high performance at a relatively low cost.
Although a cluster may consist of just a few personal computers connected by a simple network, the cluster architecture may also be used to achieve very high levels of performance. The TOP500 organization's semiannual list of the 500 fastest supercomputers often includes many clusters, e.g. the world's fastest machine in 2011 was the K computer which has a distributed memory, cluster architecture.
History
Greg Pfister has stated that clusters were not invented by any specific vendor but by customers who could not fit all their work on one computer, or needed a backup. Pfister estimates the date as some time in the 1960s. The formal engineering basis of cluster computing as a means of doing parallel work of any sort was arguably invented by Gene Amdahl of IBM, who in 1967 published what has come to be regarded as the seminal paper on parallel processing: Amdahl's Law.
The history of early computer clusters is more or less directly tied to the history of early networks, as one of the primary motivations for the development of a network was to link computing resources, creating a de facto computer cluster.
The first production system designed as a cluster was the Burroughs B5700 in the mid-1960s. This allowed up to four computers, each with either one or two processors, to be tightly coupled to a common disk storage subsystem in order to distribute the workload. Unlike standard multiprocessor systems, each computer could be restarted without disrupting overall operation.
The first commercial loosely coupled clustering product was Datapoint Corporation's "Attached Resource Computer" (ARC) system, developed in 1977, and using ARCnet as the cluster interface. Clustering per se did not really take off until Digital Equipment Corporation released their VAXcluster product in 1984 for the VMS operating system. The ARC and VAXcluster products not only supported parallel computing, but also shared file systems and peripheral devices. The idea was to provide the advantages of parallel processing, while maintaining data reliability and uniqueness. Two other noteworthy early commercial clusters were the Tandem NonStop (a 1976 high-availability commercial product) and the IBM S/390 Parallel Sysplex (circa 1994, primarily for business use).
Within the same time frame, while computer clusters used parallelism outside the computer on a commodity network, supercomputers began to use them within the same computer. Following the success of the CDC 6600 in 1964, the Cray 1 was delivered in 1976, and introduced internal parallelism via vector processing. While early supercomputers excluded clusters and relied on shared memory, in time some of the fastest supercomputers (e.g. the K computer) relied on cluster architectures.
Attributes of clusters
Computer clusters may be configured for different purposes ranging from general purpose business needs such as web-service support, to computation-intensive scientific calculations. In either case, the cluster may use a high-availability approach. Note that the attributes described below are not exclusive and a "computer cluster" may also use a high-availability approach, etc.
"Load-balancing" clusters are configurations in which cluster-nodes share computational workload to provide better overall performance. For example, a web server cluster may assign different queries to different nodes, so the overall response time will be optimized. However, approaches to load-balancing may significantly differ among applications, e.g. a high-performance cluster used for scientific computations would balance load with different algorithms from a web-server cluster which may just use a simple round-robin method by assigning each new request to a different node.
Computer clusters are used for computation-intensive purposes, rather than handling IO-oriented operations such as web service or databases. For instance, a computer cluster might support computational simulations of vehicle crashes or weather. Very tightly coupled computer clusters are designed for work that may approach "supercomputing".
"High-availability clusters" (also known as failover clusters, or HA clusters) improve the availability of the cluster approach. They operate by having redundant nodes, which are then used to provide service when system components fail. HA cluster implementations attempt to use redundancy of cluster components to eliminate single points of failure. There are commercial implementations of High-Availability clusters for many operating systems. The Linux-HA project is one commonly used free software HA package for the Linux operating system.
Benefits
Clusters are primarily designed with performance in mind, but installations are based on many other factors. Fault tolerance (the ability of a system to continue operating despite a malfunctioning node) enables scalability, and in high-performance situations, allows for a low frequency of maintenance routines, resource consolidation (e.g., RAID), and centralized management. Advantages include enabling data recovery in the event of a disaster and providing parallel data processing and high processing capacity.
In terms of scalability, clusters provide this in their ability to add nodes horizontally. This means that more computers may be added to the cluster, to improve its performance, redundancy and fault tolerance. This can be an inexpensive solution for a higher performing cluster compared to scaling up a single node in the cluster. This property of computer clusters can allow for larger computational loads to be executed by a larger number of lower performing computers.
When adding a new node to a cluster, reliability increases because the entire cluster does not need to be taken down. A single node can be taken down for maintenance, while the rest of the cluster takes on the load of that individual node.
If you have a large number of computers clustered together, this lends itself to the use of distributed file systems and RAID, both of which can increase the reliability and speed of a cluster.
Design and configuration
One of the issues in designing a cluster is how tightly coupled the individual nodes may be. For instance, a single computer job may require frequent communication among nodes: this implies that the cluster shares a dedicated network, is densely located, and probably has homogeneous nodes. The other extreme is where a computer job uses one or few nodes, and needs little or no inter-node communication, approaching grid computing.
In a Beowulf cluster, the application programs never see the computational nodes (also called slave computers) but only interact with the "Master" which is a specific computer handling the scheduling and management of the slaves. In a typical implementation the Master has two network interfaces, one that communicates with the private Beowulf network for the slaves, the other for the general purpose network of the organization. The slave computers typically have their own version of the same operating system, and local memory and disk space. However, the private slave network may also have a large and shared file server that stores global persistent data, accessed by the slaves as needed.
A special purpose 144-node DEGIMA cluster is tuned to running astrophysical N-body simulations using the Multiple-Walk parallel tree code, rather than general purpose scientific computations.
Due to the increasing computing power of each generation of game consoles, a novel use has emerged where they are repurposed into High-performance computing (HPC) clusters. Some examples of game console clusters are Sony PlayStation clusters and Microsoft Xbox clusters. Another example of consumer game product is the Nvidia Tesla Personal Supercomputer workstation, which uses multiple graphics accelerator processor chips. Besides game consoles, high-end graphics cards too can be used instead. The use of graphics cards (or rather their GPU's) to do calculations for grid computing is vastly more economical than using CPU's, despite being less precise. However, when using double-precision values, they become as precise to work with as CPU's and are still much less costly (purchase cost).
Computer clusters have historically run on separate physical computers with the same operating system. With the advent of virtualization, the cluster nodes may run on separate physical computers with different operating systems which are painted above with a virtual layer to look similar. The cluster may also be virtualized on various configurations as maintenance takes place; an example implementation is Xen as the virtualization manager with Linux-HA.
Data sharing and communication
Data sharing
As the computer clusters were appearing during the 1980s, so were supercomputers. One of the elements that distinguished the three classes at that time was that the early supercomputers relied on shared memory. Clusters do not typically use physically shared memory, while many supercomputer architectures have also abandoned it.
However, the use of a clustered file system is essential in modern computer clusters. Examples include the IBM General Parallel File System, Microsoft's Cluster Shared Volumes or the Oracle Cluster File System.
Message passing and communication
Two widely used approaches for communication between cluster nodes are MPI (Message Passing Interface) and PVM (Parallel Virtual Machine).
PVM was developed at the Oak Ridge National Laboratory around 1989 before MPI was available. PVM must be directly installed on every cluster node and provides a set of software libraries that paint the node as a "parallel virtual machine". PVM provides a run-time environment for message-passing, task and resource management, and fault notification. PVM can be used by user programs written in C, C++, or Fortran, etc.
MPI emerged in the early 1990s out of discussions among 40 organizations. The initial effort was supported by ARPA and National Science Foundation. Rather than starting anew, the design of MPI drew on various features available in commercial systems of the time. The MPI specifications then gave rise to specific implementations. MPI implementations typically use TCP/IP and socket connections. MPI is now a widely available communications model that enables parallel programs to be written in languages such as C, Fortran, Python, etc. Thus, unlike PVM which provides a concrete implementation, MPI is a specification which has been implemented in systems such as MPICH and Open MPI.
Cluster management
One of the challenges in the use of a computer cluster is the cost of administrating it which can at times be as high as the cost of administrating N independent machines, if the cluster has N nodes. In some cases this provides an advantage to shared memory architectures with lower administration costs. This has also made virtual machines popular, due to the ease of administration.
Task scheduling
When a large multi-user cluster needs to access very large amounts of data, task scheduling becomes a challenge. In a heterogeneous CPU-GPU cluster with a complex application environment, the performance of each job depends on the characteristics of the underlying cluster. Therefore, mapping tasks onto CPU cores and GPU devices provides significant challenges. This is an area of ongoing research; algorithms that combine and extend MapReduce and Hadoop have been proposed and studied.
Node failure management
When a node in a cluster fails, strategies such as "fencing" may be employed to keep the rest of the system operational. Fencing is the process of isolating a node or protecting shared resources when a node appears to be malfunctioning. There are two classes of fencing methods; one disables a node itself, and the other disallows access to resources such as shared disks.
The STONITH method stands for "Shoot The Other Node In The Head", meaning that the suspected node is disabled or powered off. For instance, power fencing uses a power controller to turn off an inoperable node.
The resources fencing approach disallows access to resources without powering off the node. This may include persistent reservation fencing via the SCSI3, fibre channel fencing to disable the fibre channel port, or global network block device (GNBD) fencing to disable access to the GNBD server.
Software development and administration
Parallel programming
Load balancing clusters such as web servers use cluster architectures to support a large number of users and typically each user request is routed to a specific node, achieving task parallelism without multi-node cooperation, given that the main goal of the system is providing rapid user access to shared data. However, "computer clusters" which perform complex computations for a small number of users need to take advantage of the parallel processing capabilities of the cluster and partition "the same computation" among several nodes.
Automatic parallelization of programs remains a technical challenge, but parallel programming models can be used to effectuate a higher degree of parallelism via the simultaneous execution of separate portions of a program on different processors.
Debugging and monitoring
Developing and debugging parallel programs on a cluster requires parallel language primitives and suitable tools such as those discussed by the High Performance Debugging Forum (HPDF) which resulted in the HPD specifications. Tools such as TotalView were then developed to debug parallel implementations on computer clusters which use Message Passing Interface (MPI) or Parallel Virtual Machine (PVM) for message passing.
The University of California, Berkeley Network of Workstations (NOW) system gathers cluster data and stores them in a database, while a system such as PARMON, developed in India, allows visually observing and managing large clusters.
Application checkpointing can be used to restore a given state of the system when a node fails during a long multi-node computation. This is essential in large clusters, given that as the number of nodes increases, so does the likelihood of node failure under heavy computational loads. Checkpointing can restore the system to a stable state so that processing can resume without needing to recompute results.
Implementations
The Linux world supports various cluster software; for application clustering, there is distcc, and MPICH. Linux Virtual Server, Linux-HA – director-based clusters that allow incoming requests for services to be distributed across multiple cluster nodes. MOSIX, LinuxPMI, Kerrighed, OpenSSI are full-blown clusters integrated into the kernel that provide for automatic process migration among homogeneous nodes. OpenSSI, openMosix and Kerrighed are single-system image implementations.
Microsoft Windows computer cluster Server 2003 based on the Windows Server platform provides pieces for high-performance computing like the job scheduler, MSMPI library and management tools.
gLite is a set of middleware technologies created by the Enabling Grids for E-sciencE (EGEE) project.
slurm is also used to schedule and manage some of the largest supercomputer clusters (see top500 list).
Other approaches
Although most computer clusters are permanent fixtures, attempts at flash mob computing have been made to build short-lived clusters for specific computations. However, larger-scale volunteer computing systems such as BOINC-based systems have had more followers.
| Technology | Computer architecture concepts | null |
18950865 | https://en.wikipedia.org/wiki/Assault%20rifle | Assault rifle | An assault rifle is a select fire rifle that uses an intermediate-rifle cartridge and a detachable magazine. Assault rifles were first put into mass production and accepted into widespread service during World War II. The first assault rifle to see major usage was the German StG 44, a development of the earlier Mkb 42. While immediately after World War II, NATO countries were equipped with battle rifles, the development of the M16 rifle during the Vietnam War prompted the adoption of assault rifles by the rest of NATO. By the end of the 20th century, assault rifles had become the standard weapon in most of the world's armies, replacing full-powered rifles and submachine guns in most roles. The two most successful modern assault rifles are the AK-47 and the M16 designs and their derivatives.
Origin of term
The term assault rifle is generally attributed to Adolf Hitler, who used the German word Sturmgewehr (which translates to "assault rifle") as the new name for the MP 43 (Maschinenpistole), subsequently known as the Sturmgewehr 44. Allied propaganda suggested that the name was chosen for propaganda purposes, but the main purpose was to differentiate the Sturmgewehr from German submachine guns such as the MP 40.
It has been suggested, however, that the Heereswaffenamt was responsible for the name Sturmgewehr, and Hitler had no input besides signing the production order. Furthermore, Hitler was initially opposed to the idea of a new infantry rifle, as Germany lacked the industrial capacity to replace the 12,000,000 Karabiner 98k rifles already in service, only changing his mind once he saw it first-hand.
The StG 44 is generally considered the first selective fire military rifle to popularize the assault rifle concept. Today, the term assault rifle is used to define firearms sharing the same basic characteristics as the StG 44.
Definition
The U.S. Army defines assault rifles as "short, compact, selective-fire weapons that fire a cartridge intermediate in power between submachine gun and rifle cartridges." In this strict definition, a firearm must have at least the following characteristics to be considered an assault rifle:
It must be capable of selective fire.
It must have an intermediate-power cartridge: more power than a pistol but less than a standard rifle or battle rifle; examples of intermediate cartridges are the 7.92×33mm Kurz, the 7.62×39mm and 5.56×45mm NATO.
Its ammunition must be supplied from a detachable box magazine.
It must have an effective range of at least .
Rifles that meet most of these criteria, but not all, are not assault rifles according to the U.S. Army's definition. For example:
Select-fire rifles such as the FN FAL, M14, and H&K G3 main battle rifles are not assault rifles; they fire full-powered rifle cartridges.
Semi-automatic-only rifles like the Colt AR-15 are not assault rifles; they do not have select-fire capabilities.
Semi-automatic-only rifles with fixed magazines like the SKS are not assault rifles; they do not have detachable box magazines and are not capable of automatic fire.
Distinction from assault weapons
In the United States, selective-fire rifles are legally defined as "machine guns", and civilian ownership of those has been tightly regulated since 1934 under the National Firearms Act and since 1986 under the Firearm Owners Protection Act. However, the term "assault rifle" is often conflated with "assault weapon", a U.S. legal category with varying definitions which includes many semi-automatic weapons. This use has been described as incorrect and a misapplication of the term.
History
WWI designs
The Fedorov Avtomat (also anglicized as Federov, ) is a select-fire infantry rifle and also one of the world's first operational automatic rifles, designed by Vladimir Grigoryevich Fyodorov in 1915 and produced in the Russian Empire and later in the Russian Soviet Federative Socialist Republic. A total of 3,200 Fedorov rifles were manufactured between 1915 and 1924 in the city of Kovrov; the vast majority of them were made after 1920. The weapon saw limited combat in World War I, but was used more substantially in the Russian Civil War and in the Winter War. Some consider it to be an "early predecessor" or "ancestor" of the modern assault rifle.
Sturmgewehr 44
The Germans were the first to pioneer the assault rifle concept during World War II, based upon research that showed that most firefights happen within and that contemporary rifles were overpowered for most small arms combat. They would soon develop a select-fire intermediate powered rifle combining the firepower of a submachine gun with the range and accuracy of a rifle.
The result was the Sturmgewehr 44, an improvement of the earlier Maschinenkarabiner 42(H), and approximately half a million Sturmgewehrs were produced by the war's end. It fired a new and revolutionary intermediate powered cartridge, the 7.92×33mm Kurz. This new cartridge was developed by shortening the standard 7.92×57mm Mauser round and giving it a lighter 125-grain bullet, which limited range but allowed for more controllable automatic fire. A smaller, lighter cartridge also allowed soldiers to carry more ammunition "to support the higher consumption rate of automatic fire".
The Sturmgewehr 44 features an inexpensive, easy-to-make, stamped steel design and a 30-round detachable box magazine. This weapon was the prototype of all successful automatic rifles. Characteristically (and unlike previous rifles) it had a straight stock with the barrel under the gas cylinder to reduce the turning moment of recoil of the rifle in the shoulder and thus help reduce the tendency of shots to climb in automatic fire. The barrel and overall length were shorter than a traditional rifle and it had a pistol grip to hold the weapon more securely in automatic fire. "The principle of this weapon—the reduction of muzzle impulse to get usable automatic fire within the actual ranges of combat—was probably the most important advance in small arms since the invention of smokeless powder."
AK-47
Like the Germans, the Soviets were influenced by experience showing that most combat engagements occur within and that their soldiers were consistently outgunned by heavily armed German troops, especially those armed with Sturmgewehr 44 assault rifles. On July 15, 1943, a Sturmgewehr was demonstrated before the People's Commissariat of Arms of the USSR. The Soviets were so impressed with the Sturmgewehr that they immediately set about developing an intermediate caliber automatic rifle of their own to replace the badly outdated Mosin–Nagant bolt-action rifles and PPSh-41 submachine guns that armed most of the Red Army.
The Soviets soon developed the 7.62×39mm M43 cartridge, which was first used in the semi-automatic SKS carbine and the RPD light machine gun. Hugo Schmeisser, the designer of the Sturmgewehr, was captured after World War II, and, likely, helped develop the AK-47 assault rifle, which would quickly replace the SKS and Mosin in Soviet service. The AK-47 was finalized, adopted and entered widespread service in the Soviet army in the early 1950s. Its firepower, ease of use, low production costs, and reliability were perfectly suited for the Red Army's new mobile warfare doctrines. In the 1960s, the Soviets introduced the RPK light machine gun, itself an AK-47 type weapon with a bipod, a stronger receiver, and a longer, heavier barrel that would eventually replace the RPD light machine gun.
The AK-47 was widely supplied or sold to nations allied with the USSR, and the blueprints were shared with several friendly nations (the People's Republic of China standing out among these with the Type 56). As a result, more AK-type weapons have been produced than all other assault rifles combined. As of 2004, "of the estimated 500 million firearms worldwide, approximately 100 million belong to the Kalashnikov family, three-quarters of which are AK-47s."
Post-StG battle rifles
The U.S. Army was influenced by combat experience with semi-automatic weapons such as the M1 Garand and M1 Carbine, which enjoyed a significant advantage over enemies armed primarily with bolt-action rifles. Although U.S. Army studies of World War II combat accounts had very similar results to that of the Germans and Soviets, the U.S. Army failed to recognize the importance of the assault rifle concept, and instead maintained its traditional views and preference for high-powered semi-automatic rifles. At the time, the U.S. Army believed that the Sturmgewehr 44 was "intended in a general way to serve the same purpose as the U.S. carbine" and was in many ways inferior to the M1 carbine, and was of "little importance".
After World War II, the United States military started looking for a single automatic rifle to replace the M1 Garand, M1/M2 Carbines, M1918 Browning Automatic Rifle, M3 "Grease Gun" and Thompson submachine gun. Early experiments with select-fire versions of the M1 Garand proved disappointing. During the Korean War, the select-fire M2 Carbine largely replaced the submachine gun in U.S. service and became the most widely used Carbine variant. Combat experience suggested that the .30 Carbine round was under-powered. American weapons designers reached the same conclusion as the German and Soviet ones: an intermediate round was necessary, and recommended a small-caliber, high-velocity cartridge.
Senior American commanders had faced fanatical enemies and experienced major logistical problems during World War II and the Korean War, and insisted that a single powerful .30 caliber cartridge be developed, that could be used by the new automatic rifle, and also by the new general-purpose machine gun (GPMG) in concurrent development. This culminated in the development of the 7.62×51mm NATO cartridge and the M14 rifle which was basically an improved select-fire M1 Garand with a 20-round magazine. The U.S. also adopted the M60 GPMG, which replaced the M1919 Browning machine gun in major combat roles. Its NATO partners adopted the FN FAL and Heckler & Koch G3 rifles, as well as the FN MAG and Rheinmetall MG3 GPMGs.
The FN FAL is a 7.62×51mm, selective fire, automatic rifle produced by the Belgian armaments manufacturer Fabrique Nationale de Herstal (FN). During the Cold War it was adopted by many North Atlantic Treaty Organization (NATO) countries, most notably with the British Commonwealth as the semi-automatic L1A1. It is one of the most widely used rifles in history, having been used by more than 90 countries. The FAL was predominantly chambered for the 7.62mm NATO round, and because of its prevalence and widespread use among the armed forces of many western nations during the Cold War, it was nicknamed "The right arm of the Free World".
The Heckler & Koch G3 is a 7.62×51mm, selective fire, automatic rifle produced by the German armament manufacturer Heckler & Koch GmbH (H&K) in collaboration with the Spanish state-owned design and development agency CETME (Centro de Estudios Técnicos de Materiales Especiales). The rifle proved successful in the export market, being adopted by the armed forces of over 60 countries. After World War II, German technicians involved in developing the Sturmgewehr 45, continued their research in France at CEAM. The StG 45 mechanism was modified by Ludwig Vorgrimler and Theodor Löffler at the Mulhouse facility between 1946 and 1949. Vorgrimler later went to work at CETME in Spain and developed the line of CETME automatic rifles based on his improved StG 45 design. Germany eventually purchased the license for the CETME design and manufactured the Heckler & Koch G3 as well as an entire line of weapons built on the same system, one of the most famous being the MP5 SMG.
M16
The first confrontations between the AK-47 and the M14 ("assault rifle" vs "battle rifle") came in the early part of the Vietnam War. Battlefield reports indicated that the M14 was uncontrollable in full-auto and that soldiers could not carry enough ammunition to maintain fire superiority over the AK-47. And, while the M2 Carbine offered a high rate of fire, it was under-powered and ultimately outclassed by the AK-47. A replacement was needed: A medium between the traditional preference for high-powered rifles such as the M14, and the lightweight firepower of the M2 Carbine.
As a result, the Army was forced to reconsider a 1957 request by General Willard G. Wyman, commander of the U.S. Continental Army Command (CONARC) to develop a .223 caliber (5.56 mm) select-fire rifle weighing when loaded with a 20-round magazine. The 5.56 mm round had to penetrate a standard U.S. helmet at and retain a velocity in excess of the speed of sound, while matching or exceeding the wounding ability of the .30 Carbine cartridge.
This request ultimately resulted in the development of a scaled-down version of the ArmaLite AR-10, called the ArmaLite AR-15 rifle. However, despite overwhelming evidence that the AR-15 could bring more firepower to bear than the M14, the Army opposed the adoption of the new rifle. In January 1963, Secretary of Defense Robert McNamara concluded that the AR-15 was the superior weapon system and ordered a halt to M14 production. At the time, the AR-15 was the only rifle available that could fulfill the requirement of a universal infantry weapon for issue to all services.
After modifications (most notably, the charging handle was re-located from under the carrying handle like it was on AR-10 to the rear of the receiver), the newly redesigned rifle was subsequently adopted as the M16 Rifle. "(The M16) was much lighter compared to the M14 it replaced, ultimately allowing soldiers to carry more ammunition. The air-cooled, gas-operated, magazine-fed assault rifle was made of steel, aluminum alloy and composite plastics, truly cutting-edge for the time. Designed with full and semi-automatic capabilities, the weapon initially did not respond well to wet and dirty conditions, sometimes even jamming in combat. After a few minor modifications, the weapon gained in popularity among troops on the battlefield."
Despite its early failures, the M16 proved to be a revolutionary design and stands as the longest continuously serving rifle in American military history. It has been adopted by many U.S. allies and the 5.56×45mm NATO cartridge has become not only the NATO standard but "the standard assault-rifle cartridge in much of the world". It also led to the development of small-caliber high-velocity service rifles by every major army in the world, including the USSR and People's Republic of China. Today, many small arms experts consider the M16 the standard by which all other assault rifles are judged.
HK33
During the 1960s other countries would follow the Americans' lead and begin to develop 5.56×45mm assault rifles, most notably Germany with the Heckler & Koch HK33. The HK33 was essentially a smaller 5.56mm version of the 7.62×51mm Heckler & Koch G3 rifle. As one of the first 5.56mm assault rifles on the market, it would go on to become one of the most widely distributed assault rifles. The HK33 featured a modular design with a wide range of accessories (telescoping butt-stocks, optics, bi-pods, etc.) that could be easily removed and arranged in a variety of configurations.
5.56mm NATO
The adoption of the M16, the H&K33, and the 5.56×45mm cartridge inspired an international trend towards relatively small-sized, lightweight, high-velocity military service cartridges that allow a soldier to carry more ammunition for the same weight compared to the larger and heavier 7.62×51mm NATO cartridge. The 5.56mm cartridge is also much easier to shoot. In 1961 marksmanship testing, the U.S. Army found that 43% of AR-15 shooters achieved Expert, while only 22% of M-14 rifle shooters did so. Also, a lower recoil impulse, allows for more controllable automatic weapons fire.
In March 1970, the U.S. recommended that all NATO forces adopt the 5.56×45mm cartridge. This shift represented a change in the philosophy of the military's long-held position about caliber size. By the middle of the 1970s, other armies were looking at assault rifle-type weapons. A NATO standardization effort soon started and tests of various rounds were carried out starting in 1977. The U.S. offered the 5.56×45mm M193 round, but there were concerns about its penetration in the face of the wider introduction of body armor. In the end the Belgian 5.56×45mm SS109 round was chosen (STANAG 4172) in October 1980. The SS109 round was based on the U.S. cartridge but included a new stronger, heavier, 62-grain bullet design, with better long-range performance and improved penetration (specifically, to consistently penetrate the side of a steel helmet at ).
Also during the 1970s, Finland, Israel, and South Africa introduced AK type assault rifles in 5.56×45mm. Sweden began the transition with trails in 1981 and full adaptation in 1986. During the 1990s, Russia developed the AK-101 in 5.56×45mm NATO for the world export market. In addition, Bulgaria, Czechoslovakia, Hungary, Poland and Yugoslavia (i.e., Serbia) have also rechambered their locally produced assault rifles to 5.56mm NATO.
AK-74
The AK-74 assault rifle was a Soviet answer to the U.S. M16. The Soviet military realized that the M16 had better range and accuracy over the AKM, and that its lighter cartridge allowed soldiers to carry more ammunition. Therefore, in 1967, the USSR issued an official requirement to replace the AKM and the 7.62×39mm cartridge. They soon began to develop the AK-74 and the 5.45×39mm cartridge. AK-74 production began in 1974, and it was unveiled in 1977, when it was carried by Soviet parachute troops during the annual Red Square parade. It would soon replace the AKM and become the standard Soviet infantry rifle. In 1979, the AK-74 saw combat for the first time in Afghanistan, where the lethality of the 5.45mm rounds led to the Mujahadeen dubbing them "poison bullets".
The adoption of the 5.56mm NATO and the Russian 5.45×39mm cartridges cemented the worldwide trend toward small caliber, high-velocity cartridges.
Compact assault rifles
Following the adoption of the M16, carbine variants were also adopted for close quarters operations. The AR-15 family of weapons served through the Vietnam War. However, these compact assault rifles had design issues, as "the barrel length was halved" to which "upset the ballistics", reducing its range and accuracy and leading "to considerable muzzle flash and blast, so that a large flash suppressor had to be fitted". "Nevertheless, as a short-range weapon it is quite adequate and thus, [despite] its caliber, [the Colt Commando] is classed as a submachine gun." Other compact assault rifles, such as the HK53, AKS-74U and the Daewoo K1, have been made and they have also been called submachine guns.
Bullpups
In 1977, Austria introduced the 5.56×45mm Steyr AUG bullpup rifle, often cited as the first successful bullpup rifle, finding service with the armed forces of over twenty countries. It was highly advanced for the 1970s, combining in the same weapon the bullpup configuration, a polymer housing, dual vertical grips, an optical sight as standard, and a modular design. Highly reliable, light, and accurate, the Steyr AUG showed clearly the potential of the bullpup layout. In 1978, France introduced the 5.56×45mm FAMAS bullpup rifle. In 1985, the British introduced the 5.56×45mm L85 bullpup rifle. In the late 1990s, Israel introduced the 5.56mm NATO Tavor TAR-21. In 1997, China adopted the QBZ-95 in the new 5.8×42mm cartridge, which they claim is superior to both the 5.56×45mm and the 5.45×39mm. By the turn of the century, the bullpup assault rifle design had achieved worldwide acceptance.
Heckler & Koch G36
The Heckler & Koch G36 is a 5.56×45mm assault rifle, designed in the early 1990s by Heckler & Koch in Germany as a replacement for the heavier G3. It was accepted into service with the in 1997, replacing the G3. The G36 is gas-operated and feeds from a 30-round detachable box magazine or 100-round C-Mag drum magazine. The G36 was made with the extensive use of lightweight, corrosion-resistant synthetic materials in its design; the receiver housing, stock, trigger group (including the fire control selector and firing mechanism parts), magazine well, handguard and carrying handle are all made of a carbon fiber-reinforced polyamide. The receiver has an integrated steel barrel trunnion (with locking recesses) and a nylon 66 steel reinforced receiver. The standard Bundeswehr versions of the G36 are equipped with a unique ZF 3×4° dual optical sight that combines a 3× magnified telescopic sight and an unmagnified reflex sight mounted on top of the telescopic sight. Widely distributed, it has been adopted by over 40 countries and prompted other nations to develop similar composite designs, such as the FX-05 Xiuhcoatl.
| Technology | Projectile weapons | null |
18952324 | https://en.wikipedia.org/wiki/Pitaya | Pitaya | A pitaya () or pitahaya () is the fruit of several cactus species indigenous to the region of southern Mexico and along the Pacific coasts of Guatemala, Costa Rica, and El Salvador. Pitaya is cultivated in East Asia, South Asia, Southeast Asia, the United States, the Caribbean, Australia, Brazil, and throughout tropical and subtropical regions of the world.
Pitaya usually refers to fruit of the genus Stenocereus, while pitahaya or dragon fruit refers to fruit of the genus Selenicereus (formerly Hylocereus), both in the family Cactaceae. The common name in English dragon fruit derives from the leather-like skin and scaly spikes on the fruit exterior. Depending on the variety, pitaya fruits may have sweet- or sour-tasting flesh that can be red, white, or yellow in color.
Vernacular names
These fruits are commonly known in English as "dragon fruit", a name used since 1963, apparently resulting from the leather-like skin and prominent scaly spikes on the fruit exterior. The fruit is often designated as "Vietnamese dragon fruit" as Vietnam is the lead exporter. The fruit may also be known as a strawberry pear.
The names pitahaya and pitaya derive from Mexico, and pitaya roja in Central America and northern South America, possibly relating to pitahaya for names of tall cacti species with flowering fruit.
Geography
Pitaya or dragon fruit is native to the region of southern Mexico and along the Pacific coasts of Guatemala, Costa Rica, and El Salvador. The dragon fruit is cultivated in East Asia, South Asia, Southeast Asia, the United States, the Caribbean, Australia, and throughout tropical and subtropical regions of the world.
Varieties
Stenocereus
Stenocereus fruit (sour pitayas) are a variety that is commonly eaten in the arid regions of the Americas. They are more sour and refreshing, with juicier flesh and a stronger taste.
The sour pitaya or pitaya agria (S. gummosus) in the Sonoran Desert has been an important food source for indigenous peoples of the Americas. The Seri people of northwestern Mexico still harvest the fruit, and call the plant ziix is ccapxl "thing whose fruit is sour".
The fruit of related species, such as S. queretaroensis and the dagger cactus or pitaya de mayo (S. griseus), are also locally important foods. The fruit of the organ pipe cactus (S. thurberi, called ool by the Seris) is the pitaya dulce "sweet pitaya".
Dragon fruit, Selenicereus
Sweet pitayas come in three types, all with leathery, slightly leafy skin:
Selenicereus undatus (Pitaya blanca or white-fleshed pitaya, also known as Hylocereus undatus) has pink-skinned fruit with white flesh. This is the most commonly seen "dragon fruit".
Selenicereus costaricensis (Pitaya roja or red-fleshed pitaya, also known as Hylocereus costaricensis, and possibly incorrectly as Hylocereus polyrhizus) has red-skinned fruit with red flesh.
Selenicereus megalanthus (Pitaya amarilla or yellow pitaya, also known as Hylocereus megalanthus) has yellow-skinned fruit with white flesh.
The fruit normally weighs from ; some may reach . Early imports from Colombia to Australia were designated "Hylocereus ocampensis" (or "Cereus repandus", the red fruit) and "Cereus triangularis" (supposedly, the yellow fruit or the three-sided cross-section of the stem).
Cultivation
After a thorough cleaning of the seeds from the pulp of the fruit, the seeds may be stored when dried. The ideal fruit is unblemished and overripe.
Seeds grow well in a compost or potting soil mix – even as a potted indoor plant. Pitaya cacti usually germinate after between 11 and 14 days after shallow planting. As they are cacti, overwatering is a concern for home growers. As their growth continues, these climbing plants will find something to climb on, which can involve putting aerial roots down from the branches in addition to the basal roots. Once the plant reaches a mature in weight, the plant may flower.
Commercial plantings can be done at high density with between . Plants can take up to 60 months/260 weeks to come into full commercial production, at which stage yields of can be expected.
Pitaya flowers bloom overnight and usually wilt by the evening. They rely on nocturnal pollinators such as bats or moths for fertilization. Self-fertilization will not produce fruit in some species and while crossbreeding has resulted in several "self-fertile" varieties, cross-pollinating with a second, genetically distinct plant of the same species generally increases fruit set and quality. This limits the capability of home growers to produce the fruit. However, the plants can flower between three and six times per year depending on growing conditions. Like other cacti, if a healthy piece of the stem is broken off, it may take root in the soil and become its own plant.
The plants can endure temperatures up to and short periods of frost but will not survive long exposure to freezing temperatures. The cacti thrive most in USDA zones 10–11 but may survive outdoors in zone 9a or 9b.
Selenicereus has adapted to live in dry tropical climates with a moderate amount of rain. In numerous regions, it has escaped cultivation to become a weed and is classified as an invasive weed in some countries.
Pests and diseases
Stems and fruits are susceptible to several diseases caused by fungi, bacteria, a nematode, and a virus. Overwatering or excessive rainfall can cause the flowers to drop and fruit to rot. The bacterium Xanthomonas campestris causes the stems to rot. Dothiorella fungi can cause brown spots on the fruit. Other fungi known to infect pitaya include Botryosphaeria dothidea, Colletotrichum gloeosporioides and Bipolaris cactivora.
Uses
Culinary
The fruit's texture is sometimes likened to that of the kiwifruit because of its black, crunchy seeds. The seed oil contains the fatty acids linoleic acid and linolenic acid. Dragon fruit is used to flavor and color juices and alcoholic beverages, such as "Dragon's Blood Punch" and the "Dragotini". The flowers can be eaten or steeped as tea.
The red and purple colors of some Selenicereus fruits are due to betacyanins, a family of pigments that includes betanin, the same substance that gives beets, Swiss chard, and amaranth their red color.
Nutrients
The USDA FoodData Central database published their analysis of the nutritional contents of raw Pitaya in 2022. The majority of the fruit by weight is water (87g out of 100g). One serving of provides of food energy.
The USDA also reports one limited product label entry from a manufacturer of a branded product, showing that a reference serving of dried pitaya provides of food energy, 82% carbohydrates, 4% protein, and 11% of the Daily Value each for vitamin C and calcium.
Seed oils
The fatty acid compositions of the seed oils of Selenicereus costaricensis, syn. Hylocereus costaricensis (red-fleshed pitaya) and Selenicereus undatus, syn. Hylocereus undatus (white-fleshed pitaya) were similar: myristic acid (negligible), palmitic acid (17%), stearic acid (5%), palmitoleic acid (about 1%), oleic acid (22%), cis-vaccenic acid (3%), linoleic acid (50%), and α-linolenic acid (1%).
Gallery
| Biology and health sciences | Other culinary fruits | Plants |
18952443 | https://en.wikipedia.org/wiki/Osmoregulation | Osmoregulation | Osmoregulation is the active regulation of the osmotic pressure of an organism's body fluids, detected by osmoreceptors, to maintain the homeostasis of the organism's water content; that is, it maintains the fluid balance and the concentration of electrolytes (salts in solution which in this case is represented by body fluid) to keep the body fluids from becoming too diluted or concentrated. Osmotic pressure is a measure of the tendency of water to move into one solution from another by osmosis. The higher the osmotic pressure of a solution, the more water tends to move into it. Pressure must be exerted on the hypertonic side of a selectively permeable membrane to prevent diffusion of water by osmosis from the side containing pure water.
Although there may be hourly and daily variations in osmotic balance, an animal is generally in an osmotic steady state over the long term. Organisms in aquatic and terrestrial environments must maintain the right concentration of solutes and amount of water in their body fluids; this involves excretion (getting rid of metabolic nitrogen wastes and other substances such as hormones that would be toxic if allowed to accumulate in the blood) through organs such as the skin and the kidneys.
Regulators and conformers
Two major types of osmoregulation are osmoconformers and osmoregulators. Osmoconformers match their body osmolarity to their environment actively or passively. Most marine invertebrates are osmoconformers, although their ionic composition may be different from that of seawater. In a strictly osmoregulating animal, the amounts of internal salt and water are held relatively constant in the face of environmental changes. It requires that intake and outflow of water and salts be equal over an extended period of time.
Organisms that maintain an internal osmolarity different from the medium in which they are immersed have been termed osmoregulators. They tightly regulate their body osmolarity, maintaining constant internal conditions. They are more common in the animal kingdom. Osmoregulators actively control salt concentrations despite the salt concentrations in the environment. An example is freshwater fish. The gills actively uptake salt from the environment by the use of mitochondria-rich cells. Water will diffuse into the fish, so it excretes a very hypotonic (dilute) urine to expel all the excess water. A marine fish has an internal osmotic concentration lower than that of the surrounding seawater, so it tends to lose water and gain salt. It actively excretes salt out from the gills. Most fish are stenohaline, which means they are restricted to either salt or fresh water and cannot survive in water with a different salt concentration than they are adapted to. However, some fish show an ability to effectively osmoregulate across a broad range of salinities; fish with this ability are known as euryhaline species, e.g., flounder. Flounder have been observed to inhabit two disparate environments—marine and fresh water—and it is inherent to adapt to both by bringing in behavioral and physiological modifications.
Some marine fish, like sharks, have adopted a different, efficient mechanism to conserve water, i.e., osmoregulation. They retain urea in their blood in relatively higher concentration. Urea damages living tissues so, to cope with this problem, some fish retain trimethylamine oxide, which helps to counteract urea's destabilizing effects on cells. Sharks, having slightly higher solute concentration (i.e., above 1000 mOsm which is sea solute concentration), do not drink water like fresh water fish.
In plants
While there are no specific osmoregulatory organs in higher plants, the stomata are important in regulating water loss through evapotranspiration, and on the cellular level the vacuole is crucial in regulating the concentration of solutes in the cytoplasm. Strong winds, low humidity and high temperatures all increase evapotranspiration from leaves. Abscisic acid is an important hormone in helping plants to conserve water—it causes stomata to close and stimulates root growth so that more water can be absorbed.
Plants share with animals the problems of obtaining water but, unlike in animals, the loss of water in plants is crucial to create a driving force to move nutrients from the soil to tissues. Certain plants have evolved methods of water conservation.
Xerophytes are plants that can survive in dry habitats, such as deserts, and are able to withstand prolonged periods of water shortage. Succulent plants such as the cacti store water in the vacuoles of large parenchyma tissues. Other plants have leaf modifications to reduce water loss, such as needle-shaped leaves, sunken stomata, and thick, waxy cuticles as in the pine. The sand-dune marram grass has rolled leaves with stomata on the inner surface.
Hydrophytes are plants that grow in aquatic habitats; they may be floating, submerged, or emergent, and may grow in seasonal (rather than permanent) wetlands. In these plants the water absorption may occur through the whole surface of the plant, e.g., the water lily, or solely through the roots, as in sedges. These plants do not face major osmoregulatory challenges from water scarcity, but aside from species adapted for seasonal wetlands, have few defenses against desiccation.
Halophytes are plants living in soils with high salt concentrations, such as salt marshes or alkaline soils in desert basins. They have to absorb water from such a soil which has higher salt concentration and therefore lower water potential(higher osmotic pressure). Halophytes cope with this situation by activating salts in their roots. As a consequence, the cells of the roots develop lower water potential which brings in water by osmosis. The excess salt can be stored in cells or excreted out from salt glands on leaves. The salt thus secreted by some species help them to trap water vapours from the air, which is absorbed in liquid by leaf cells. Therefore, this is another way of obtaining additional water from air, e.g., glasswort and cord-grass.
Mesophytes are plants living in lands of temperate zone, which grow in well-watered soil. They can easily compensate the water lost by transpiration through absorbing water from the soil. To prevent excessive transpiration they have developed a waterproof external covering called cuticle.
In animals
Humans
Kidneys play a very large role in human osmoregulation by regulating the amount of water reabsorbed from glomerular filtrate in kidney tubules, which is controlled by hormones such as antidiuretic hormone (ADH), aldosterone, and angiotensin II. For example, a decrease in water potential is detected by osmoreceptors in the hypothalamus, which stimulates ADH release from the pituitary gland to increase the permeability of the walls of the collecting ducts in the kidneys. Therefore, a large proportion of water is reabsorbed from fluid in the kidneys to prevent too much water from being excreted.
Marine mammals
Drinking is not common behavior in pinnipeds and cetaceans. Water balance is maintained in marine mammals by metabolic and dietary water, while accidental ingestion and dietary salt may help maintain homeostasis of electrolytes. The kidneys of pinnipeds and cetaceans are lobed in structure, unlike those of non-bears among terrestrial mammals, but this specific adaptation does not confer any greater concentrating ability. Unlike most other aquatic mammals, manatees frequently drink fresh water and sea otters frequently drink saltwater.
Teleosts
In teleost (advanced ray-finned) fishes, the gills, kidney and digestive tract are involved in maintenance of body fluid balance, as the main osmoregulatory organs. Gills in particular are considered the primary organ by which ionic concentration is controlled in marine teleosts.
Unusually, the catfishes in the eeltail family Plotosidae have an extra-branchial salt-secreting dendritic organ. The dendritic organ is likely a product of convergent evolution with other vertebrate salt-secreting organs. The role of this organ was discovered by its high NKA and NKCC activity in response to increasing salinity. However, the Plotosidae dendritic organ may be of limited use under extreme salinity conditions, compared to more typical gill-based ionoregulation.
In protists
Amoeba makes use of contractile vacuoles to collect excretory wastes, such as ammonia, from the intracellular fluid by diffusion and active transport. As osmotic action pushes water from the environment into the cytoplasm, the vacuole moves to the surface and pumps the contents into the environment.
In bacteria
Bacteria respond to osmotic stress by rapidly accumulating electrolytes or small organic solutes via transporters whose activities are stimulated by increases in osmolarity. The bacteria may also turn on genes encoding transporters of osmolytes and enzymes that synthesize osmoprotectants. The EnvZ/OmpR two-component system, which regulates the expression of porins, is well characterized in the model organism E. coli.
Vertebrate excretory systems
Waste products of the nitrogen metabolism
Ammonia is a toxic by-product of protein metabolism and is generally converted to less toxic substances after it is produced then excreted; mammals convert ammonia to urea, whereas birds and reptiles form uric acid to be excreted with other wastes via their cloacas.
Achieving osmoregulation in vertebrates
Four processes occur:
filtration – fluid portion of blood (plasma) is filtered from a nephron (functional unit of vertebrate kidney) structure known as the glomerulus into Bowman's capsule or glomerular capsule (in the kidney's cortex) and flows down the proximal convoluted tubule to a "u-turn" called the Loop of Henle (loop of the nephron) in the medulla portion of the kidney.
reabsorption – most of the viscous glomerular filtrate is returned to blood vessels that surround the convoluted tubules.
secretion – the remaining fluid becomes urine, which travels down collecting ducts to the medullary region of the kidney.
excretion – the urine (in mammals) is stored in the urinary bladder and exits via the urethra; in other vertebrates, the urine mixes with other wastes in the cloaca before leaving the body (frogs also have a urinary bladder).
| Biology and health sciences | Basics | Biology |
18952492 | https://en.wikipedia.org/wiki/Anthocyanin | Anthocyanin | Anthocyanins (), also called anthocyans, are water-soluble vacuolar pigments that, depending on their pH, may appear red, purple, blue, or black. In 1835, the German pharmacist Ludwig Clamor Marquart named a chemical compound that gives flowers a blue color, Anthokyan, in his treatise "Die Farben der Blüthen" (English: The Colors of Flowers). Food plants rich in anthocyanins include the blueberry, raspberry, black rice, and black soybean, among many others that are red, blue, purple, or black. Some of the colors of autumn leaves are derived from anthocyanins.
Anthocyanins belong to a parent class of molecules called flavonoids synthesized via the phenylpropanoid pathway. They can occur in all tissues of higher plants, including leaves, stems, roots, flowers, and fruits. Anthocyanins are derived from anthocyanidins by adding sugars. They are odorless and moderately astringent.
Although approved as food and beverage colorant in the European Union, anthocyanins are not approved for use as a food additive because they have not been verified as safe when used as food or supplement ingredients. There is no conclusive evidence that anthocyanins have any effect on human biology or diseases.
Anthocyanin-rich plants
Coloration
In flowers, the coloration that is provided by anthocyanin accumulation may attract a wide variety of animal pollinators, while in fruits, the same coloration may aid in seed dispersal by attracting herbivorous animals to the potentially-edible fruits bearing these red, blue, or purple colors.
Plant physiology
Anthocyanins may have a protective role in plants against extreme temperatures. Tomato plants protect against cold stress with anthocyanins countering reactive oxygen species, leading to a lower rate of cell death in leaves.
Light absorbance
The absorbance pattern responsible for the red color of anthocyanins may be complementary to that of green chlorophyll in photosynthetically active tissues such as young Quercus coccifera leaves. It may protect the leaves from attacks by herbivores that may be attracted by green color.
Occurrence
Anthocyanins are found in the cell vacuole, mostly in flowers and fruits, but also in leaves, stems, and roots. In these parts, they are found predominantly in outer cell layers such as the epidermis and peripheral mesophyll cells.
Most frequently occurring in nature are the glycosides of cyanidin, delphinidin, malvidin, pelargonidin, peonidin, and petunidin. Roughly 2% of all hydrocarbons fixed in photosynthesis are converted into flavonoids and their derivatives, such as the anthocyanins. Not all land plants contain anthocyanin; in the Caryophyllales (including cactus, beets, and amaranth), they are replaced by betalains. Anthocyanins and betalains have never been found in the same plant.
Sometimes bred purposely for high anthocyanin content, ornamental plants such as sweet peppers may have unusual culinary and aesthetic appeal.
In flowers
Anthocyanins occur in the flowers of many plants, such as the blue poppies of some Meconopsis species and cultivars. Anthocyanins have also been found in various tulip flowers, such as Tulipa gesneriana, Tulipa fosteriana and Tulipa eichleri.
In food
Plants rich in anthocyanins are Vaccinium species, such as blueberry, cranberry, and bilberry; Rubus berries, including black raspberry, red raspberry, and blackberry; blackcurrant, cherry, eggplant (aubergine) peel, black rice, ube, Okinawan sweet potato, Concord grape, muscadine grape, red cabbage, and violet petals. Red-fleshed peaches and apples contain anthocyanins. Anthocyanins are less abundant in banana, asparagus, pea, fennel, pear, and potato, and may be totally absent in certain cultivars of green gooseberries.
The highest recorded amount appears to be specifically in the seed coat of black soybean (Glycine max L. Merr.) containing approximately 2 g per 100 g, in purple corn kernels and husks, and in the skins and pulp of black chokeberry (Aronia melanocarpa L.) (see table). Due to critical differences in sample origin, preparation, and extraction methods determining anthocyanin content, the values presented in the adjoining table are not directly comparable.
Nature, traditional agriculture methods, and plant breeding have produced various uncommon crops containing anthocyanins, including blue- or red-flesh potatoes and purple or red broccoli, cabbage, cauliflower, carrots, and corn. Garden tomatoes have been subjected to a breeding program using introgression lines of genetically modified organisms (but not incorporating them in the final purple tomato) to define the genetic basis of purple coloration in wild species that originally were from Chile and the Galapagos Islands. The variety known as "Indigo Rose" became available commercially to the agricultural industry and home gardeners in 2012. Investing tomatoes with high anthocyanin content doubles their shelf-life and inhibits growth of a post-harvest mold pathogen, Botrytis cinerea.
Some tomatoes also have been modified genetically with transcription factors from snapdragons to produce high levels of anthocyanins in the fruits. Anthocyanins also may be found in naturally ripened olives, and are partly responsible for the red and purple colors of some olives.
In leaves of plant foods
Content of anthocyanins in the leaves of colorful plant foods such as purple corn, blueberries, or lingonberries, is about ten times higher than in the edible kernels or fruit.
The color spectrum of grape berry leaves may be analysed to evaluate the amount of anthocyanins. Fruit maturity, quality, and harvest time may be evaluated on the basis of the spectrum analysis.
Autumn leaf color
The reds, purples, and their blended combinations responsible for autumn foliage are derived from anthocyanins. Unlike carotenoids, anthocyanins are not present in the leaf throughout the growing season, but are produced actively, toward the end of summer. They develop in late summer in the sap of leaf cells, resulting from complex interactions of factors inside and outside the plant. Their formation depends on the breakdown of sugars in the presence of light as the level of phosphate in the leaf is reduced. Orange leaves in autumn result from a combination of anthocyanins and carotenoids.
Anthocyanins are present in approximately 10% of tree species in temperate regions, although in certain areas such as New England, up to 70% of tree species may produce anthocyanins.
Colorant safety
Anthocyanins are approved for use as food colorants in the European Union, Australia, and New Zealand, having colorant code E163. In 2013, a panel of scientific experts for the European Food Safety Authority concluded that anthocyanins from various fruits and vegetables have been insufficiently characterized by safety and toxicology studies to approve their use as food additives. Extending from a safe history of using red grape skin extract and blackcurrant extracts to color foods produced in Europe, the panel concluded that these extract sources were exceptions to the ruling and were sufficiently shown to be safe.
Anthocyanin extracts are not specifically listed among approved color additives for foods in the United States; however, grape juice, red grape skin and many fruit and vegetable juices, which are approved for use as colorants, are rich in naturally occurring anthocyanins. No anthocyanin sources are included among approved colorants for drugs or cosmetics. When esterified with fatty acids, anthocyanins can be used as a lipophilic colorant for foods.
In human consumption
Although anthocyanins have been shown to have antioxidant properties in vitro, there is no evidence for antioxidant effects in humans after consuming foods rich in anthocyanins. Unlike controlled test-tube conditions, the fate of anthocyanins in vivo shows they are poorly conserved (less than 5%), with most of what is absorbed existing as chemically modified metabolites that are excreted rapidly. The increase in antioxidant capacity of blood seen after the consumption of anthocyanin-rich foods may not be caused directly by the anthocyanins in the food, but instead by increased uric acid levels derived from metabolizing flavonoids (anthocyanin parent compounds) in the food. It is possible that metabolites of ingested anthocyanins are reabsorbed in the gastrointestinal tract from where they may enter the blood for systemic distribution and have effects as smaller molecules.
In a 2010 review of scientific evidence concerning the possible health benefits of eating foods claimed to have "antioxidant properties" due to anthocyanins, the European Food Safety Authority concluded that 1) there was no basis for a beneficial antioxidant effect from dietary anthocyanins in humans, 2) there was no evidence of a cause-and-effect relationship between the consumption of anthocyanin-rich foods and protection of DNA, proteins, and lipids from oxidative damage, and 3) there was no evidence generally for consumption of anthocyanin-rich foods having any "antioxidant", "anti-cancer", "anti-aging", or "healthy aging" effects.
Chemical properties
Flavylium cation derivatives
Glycosides of anthocyanidins
The anthocyanins, anthocyanidins with sugar group(s), are mostly 3-glucosides of the anthocyanidins. The anthocyanins are subdivided into the sugar-free anthocyanidin aglycones and the anthocyanin glycosides. As of 2003, more than 400 anthocyanins had been reported, while later literature in early 2006, puts the number at more than 550 different anthocyanins. The difference in chemical structure that occurs in response to changes in pH, is the reason why anthocyanins often are used as pH indicators, as they change from red in acids to blue in bases through a process called halochromism.
Stability
Anthocyanins are thought to be subject to physiochemical degradation in vivo and in vitro. Structure, pH, temperature, light, oxygen, metal ions, intramolecular association, and intermolecular association with other compounds (copigments, sugars, proteins, degradation products, etc.) generally are known to affect the color and stability of anthocyanins. B-ring hydroxylation status and pH have been shown to mediate the degradation of anthocyanins to their phenolic acid and aldehyde constituents. Indeed, significant portions of ingested anthocyanins are likely to degrade to phenolic acids and aldehyde in vivo, following consumption. This characteristic confounds scientific isolation of specific anthocyanin mechanisms in vivo.
pH
Anthocyanins generally are degraded at higher pH. However, some anthocyanins, such as petanin (petunidin 3-[6-O-(4-O-(E)-p-coumaroyl-O-α--rhamnopyranosyl)-β--glucopyranoside]-5-O-β--glucopyranoside), are resistant to degradation at pH 8 and may be used effectively as a food colorant.
Use as environmental pH indicator
Anthocyanins may be used as pH indicators because their color changes with pH; they are red or pink in acidic solutions (pH < 7), purple in neutral solutions (pH ≈ 7), greenish-yellow in alkaline solutions (pH > 7), and colorless in very alkaline solutions, where the pigment is completely reduced.
Biosynthesis
Anthocyanin pigments are assembled like all other flavonoids from two different streams of chemical raw materials in the cell:
One stream involves the shikimate pathway to produce the amino acid phenylalanine, (see phenylpropanoids)
The other stream produces three molecules of malonyl-CoA, a C3 unit from a C2 unit (acetyl-CoA),
These streams meet and are coupled together by the enzyme chalcone synthase, which forms an intermediate chalcone-like compound via a polyketide folding mechanism that is commonly found in plants,
The chalcone is subsequently isomerized by the enzyme chalcone isomerase to the prototype pigment naringenin,
Naringenin is subsequently oxidized by enzymes such as flavanone hydroxylase, flavonoid 3'-hydroxylase, and flavonoid 3',5'-hydroxylase,
These oxidation products are further reduced by the enzyme dihydroflavonol 4-reductase to the corresponding colorless leucoanthocyanidins,
Leucoanthocyanidins once were believed to be the immediate precursors of the next enzyme, a dioxygenase referred to as anthocyanidin synthase, or, leucoanthocyanidin dioxygenase. Flavan-3-ols, the products of leucoanthocyanidin reductase (LAR), recently have been shown to be their true substrates,
The resulting unstable anthocyanidins are further coupled to sugar molecules by enzymes such as UDP-3-O-glucosyltransferase, to yield the final relatively-stable anthocyanins.
Thus, more than five enzymes are required to synthesize these pigments, each working in concert. Even a minor disruption in any of the mechanisms of these enzymes by either genetic or environmental factors, would halt anthocyanin production. While the biological burden of producing anthocyanins is relatively high, plants benefit significantly from the environmental adaptation, disease tolerance, and pest tolerance provided by anthocyanins.
In anthocyanin biosynthetic pathway, L-phenylalanine is converted to naringenin by phenylalanine ammonialyase, cinnamate 4-hydroxylase, 4-coumarate CoA ligase, chalcone synthase, and chalcone isomerase. Then, the next pathway is catalyzed, resulting in the formation of complex aglycone and anthocyanin through composition by flavanone 3-hydroxylase, flavonoid 3'-hydroxylase, dihydroflavonol 4-reductase, anthocyanidin synthase, UDP-glucoside: flavonoid glucosyltransferase, and methyl transferase.
Genetic analysis
The phenolic metabolic pathways and enzymes may be studied by mean of transgenesis of genes. The Arabidopsis regulatory gene in the production of anthocyanin pigment 1 (AtPAP1) may be expressed in other plant species.
Dye-sensitized solar cells
Anthocyanins have been used in organic solar cells because of their ability to convert light energy into electrical energy. The many benefits to using dye-sensitized solar cells instead of traditional p-n junction silicon cells, include lower purity requirements and abundance of component materials, as well as the fact that they may be produced on flexible substrates, making them amenable to roll-to-roll printing processes.
Visual markers
Anthocyanins fluoresce, enabling a tool for plant cell research to allow live cell imaging without a requirement for other fluorophores. Anthocyanin production may be engineered into genetically modified materials to enable their identification visually.
| Physical sciences | Polyphenols | Chemistry |
18952520 | https://en.wikipedia.org/wiki/Euphorbiaceae | Euphorbiaceae | Euphorbiaceae (), the spurge family, is a large family of flowering plants. In English, they are also commonly called euphorbias, which is also the name of the type genus of the family. Most spurges, such as Euphorbia paralias, are herbs, but some, especially in the tropics, are shrubs or trees, such as Hevea brasiliensis. Some, such as Euphorbia canariensis, are succulent and resemble cacti because of convergent evolution. This family has a cosmopolitan global distribution. The greatest diversity of species is in the tropics; however, the Euphorbiaceae also have many species in nontropical areas of all continents except Antarctica.
Description
The leaves are alternate, seldom opposite, with stipules. They are mainly simple, but where compound, are always palmate, never pinnate. Stipules may be reduced to hairs, glands, or spines, or in succulent species are sometimes absent.
The plants can be monoecious or dioecious. The radially symmetrical flowers are unisexual, with the male and female flowers usually on the same plant. As can be expected from such a large family, a wide variety exists in the structure of the flowers. The stamens (the male organs) number from one to 10 (or even more). The female flowers are hypogynous, that is, with superior ovaries.
The genera in tribe Euphorbieae, subtribe Euphorbiinae (Euphorbia and close relatives) show a highly specialized form of pseudanthium ("false flower" made up of several true flowers) called a cyathium. This is usually a small, cup-like involucre consisting of fused-together bracts and peripheral nectary glands, surrounding a ring of male flowers, each a single stamen. In the middle of the cyathium stands a female flower, a single pistil with branched stigmas. This whole arrangement resembles a single flower.
The fruit is usually a schizocarp, but sometimes a drupe. A typical schizocarp is the regma, a capsular fruit with three or more cells, each of which splits open explosively at maturity, scattering the small seeds.
The family contains a large variety of phytotoxins (toxic substances produced by plants), including diterpene esters, alkaloids, and cyanogenic glycosides (e.g. root tubers of cassava). The seeds of the castor oil plant Ricinus communis contain the highly toxic carbohydrate-binding protein ricin.
A milky latex is a characteristic of the subfamilies Euphorbioideae and Crotonoideae, and the latex of the rubber tree Hevea brasiliensis is the primary source of natural rubber. The latex is poisonous in the Euphorbioideae, but innocuous in the Crotonoideae. White mangrove, also known as blind-your-eye mangrove latex (Excoecaria agallocha), causes blistering on contact and temporary blindness if it contacts the eyes, hence its name. The latex of spurge was used as a laxative.
Twenty first century molecular studies have shown that the enigmatic family Rafflesiaceae, which was only recently recognized to belong to order Malpighiales, is derived from within the Euphorbiaceae.
Euphorbiaceae are monoecious and open pollinated and so self-incompatibility is rare - although it has been reported in the past, apparently this was in error. It is confirmed to be absent or incomplete in herbaceous Chamaesyce by Ehrenfeld 1976, Hevea by Bouharmont 1962, and Manihot by Jennings 1963 and George & Shifriss 1967.
Taxonomy
The family Euphorbiaceae is the fifth-largest flowering plant family and has about 7,500 species organised into 300 genera, 37 tribes, and three subfamilies: Acalyphoideae, Crotonoideae and Euphorbioideae. Amongst the oldest fossils of the group include the permineralised fruit Euphorbiotheca deccanensis from the Intertrappean Beds of India, dating to the late Maastrichtian at the end of the Cretaceous, around 66 million years ago.
Uses and toxicity
Some species of Euphorbiaceae have economic significance, such as cassava (Manihot esculenta), castor oil plant (Ricinus communis), Barbados nut (Jatropha curcas), and the Pará rubber tree (Hevea brasiliensis). Many are grown as ornamental plants, such as poinsettia (Euphorbia pulcherrima) or garden croton (Codiaeum variegatum). Leafy spurge (Euphorbia esula) and Chinese tallow (Triadica sebifera) are invasive weeds in North America.
Seeds of the castor oil plant (Ricinus communis L.) contain the extremely potent toxin, ricin.
Although some species of the Euphorbiaceae have been used in traditional medicine, , there is no rigorous clinical evidence that euphorbia extracts are effective for treating any disease.
There is evidence that euphol, a tetracyclic triterpene alcohol, and the main constituent of the sap of the medicinal plant Euphorbia tirucalli, has anti-cancer activity. Analysis of toxicological screening of the inhibitory effect and bioactivity of euphol has shown concentration-dependent cytotoxic effects on cancer cell lines, with more than a five-fold difference in the IC50 values in some cell lines. Euphol treatment had a higher selective cytotoxicity index (0.64-3.36) than temozolomide (0.11-1.13) and reduced both proliferation and cell motility. Euphol also exhibited antitumoral and antiangiogenic activity in vivo, using the chicken chorioallantoic membrane assay, with synergistic temozolomide interactions in most cell lines. In conclusion, euphol exerted in vitro and in vivo cytotoxicity against glioma cells, through several cancer pathways, including the activation of autophagy-associated cell death.
Numerous Euphorbiaceae species are listed on the poisonous plant database of the US Food and Drug Administration mainly because of the toxic sap.
Phytochemistry
Phytochemicals found in Euphorbiaceae species include diterpenoids, terpenoids, flavonoids, alkaloids, tannins, neriifolins (also found in oleander), cycloartenol, lectin, and taraxerol, among others.
Conservation
Some species of this family are facing the risk of extinction. These include the Euphorbia species E. appariciana, E. attastoma, E. crossadenia, and E. gymnoclada.
| Biology and health sciences | Malpighiales | null |
18952649 | https://en.wikipedia.org/wiki/Pronghorn | Pronghorn | The pronghorn (, ) (Antilocapra americana) is a species of artiodactyl (even-toed, hoofed) mammal indigenous to interior western and central North America. Though not an antelope, it is known colloquially in North America as the American antelope, prong buck, pronghorn antelope and prairie antelope, because it closely resembles the antelopes of the Old World and fills a similar ecological niche due to parallel evolution. It is the only surviving member of the family Antilocapridae.
During the Pleistocene epoch, about 11 other antilocaprid species existed in North America, many with long or spectacularly-twisted horns. Three other genera (Capromeryx, Stockoceros and Tetrameryx) existed when humans entered North America but are now extinct.
The pronghorn's closest living relatives are the giraffe and okapi. The Antilocaprids are part of the infraorder Pecora, making them distant relatives of deer, bovids, and moschids.
The pronghorn is the fastest land mammal in the Americas, with running speeds of up to . It is the symbol of the American Society of Mammalogists.
Etymology
The animal gets its name from its horn sheaths that branch and have a forward-pointing tine, unlike the horns of species from the ox family Bovidae.
European discovery
Pronghorns were first seen and described by Spanish explorers in the 16th century, but the species was not formally recorded or scrutinized until the expedition in 1804–06 by Captain Meriwether Lewis and Second Lieutenant William Clark. Following the discovery of a few subspecies of the sharp-tailed grouse, Lewis and Clark came across the pronghorn near the mouth of the Niobrara River, in present-day Nebraska. Clark was among the first Euro-Americans to publish the experience of killing a pronghorn, and described his experience as follows:
The pronghorn was first officially described by American ornithologist George Ord in 1815.
Description
Pronghorns have distinct white fur on their rumps, sides, breasts, bellies, and across their throats. Adult males are long from nose to tail, stand high at the shoulder, and weigh . The females are the same height as males, but weigh . The feet have two hooves, with no dewclaws. Their body temperature is .
Head
They have very large eyes with a 320° field of vision. Their orbits (eye sockets) are prominent and set high on the skull. Their teeth are hypsodont, and their dental formula is . Unlike deer, pronghorns possess a gallbladder.
Each horn of the pronghorn is composed of a slender, laterally flattened blade of bone which is thought to grow from the frontal bones of the skull, or from the subcutaneous tissues of the scalp, forming a permanent core. As in the Giraffidae, skin covers the bony cores, but in the pronghorn, it develops into a keratinous sheath which is shed and regrown annually. Males have a horn sheath about (average ) long with a prong. Females have smaller horns that range from (average ) and sometimes barely visible; they are straight and very rarely pronged.
Males are further differentiated from females in having a small patch of black hair at the angle of the mandible. Pronghorns have a distinct, musky odor. Males mark territory with a preorbital scent gland which is on the sides of the head.
Scent Glands
The preorbital gland's secretion contains the highly odoriferous compound, 2-ethyl-3-methylpyrazine. This compound is also the major volatile component found on the animal's back in the male's medial gland.
Male and female animals have glands that are exposed when the white hair on the rump stands up. 2-Pyrrolidinone, the major compound in the rump gland has an odor reminiscent of buttered popcorn to humans. The flared rump hair and odor alert adjacent animals of a possible danger.
Pronghorns have well developed glands on each hoof. Like many ungulates, these interdigital (hoof) glands of pronghorn contain chemical compounds that are known to have antimicrobial activity against soil and mammalian pathogens.
Movement
The pronghorn is the fastest land mammal in the Western Hemisphere, being built for maximum predator evasion through running. The top speed is dependent upon the length of time over which it is measured. It can run for , for , and for . Although it is slower than the African cheetah, it can sustain top speeds much longer than cheetahs. The pronghorn may have evolved its running ability to escape from now-extinct predators such as the American cheetah, since its speed greatly exceeds that of all extant North American predators. Carbon and nitrogen isotope comparisons between pronghorn, horses, Bighorn sheep, bison, American cheetahs, American lions, and wolves of the Natural Trap Cave found that while American cheetahs seemed to subsist on pronghorns, they did not do so exclusively. In fact, pronghorns were also important prey of American lions and wolves. Compared to its body size, the pronghorn has a large windpipe, heart, blood volume, erythrocites and lungs to allow it to take in large amounts of air when running. Additionally, pronghorn hooves have two long, cushioned, pointed toes which help absorb shock when running at high speeds. They also have an extremely light bone structure and hollow hair.
Male pronghorn tend to have a higher level of physical activity than females and apparently also have a greater blood volume relative to body size.
Pronghorns are built for speed, not for jumping. Since their ranges are sometimes affected by sheep ranchers' fences, they can be seen going under fences, sometimes at high speed. For this reason, the Arizona Antelope Foundation and others are in the process of removing the bottom barbed wire from the fences, and/or installing a barbless bottom wire.
The pronghorn has been observed to have at least 13 distinct gaits, including one reaching nearly per stride.
When a pronghorn sees something that alarms it, the white hair on the rump flairs open and exposes two highly odoriferous glands that releases a compound described as having an odour "reminiscent of buttered popcorn." This sends a message to other pronghorns by both sight and smell about a present danger. This scent has been observed by humans 20 to 30 meters downwind from alarmed animals. The major odour compound identified from this gland is 2-pyrrolidinone.
Range and ecology
Prior to the arrival of Europeans, the pronghorn was particularly abundant in the regions west of the Mississippi River (still its primary range today). Pronghorn herds filled a vital ecological niche of the prairie habitat, as well as other climatic zones. The amount of wildlife was considered to be so vast at one time that the prehistoric American Prairie—and as recently as 200 to 300 years ago—has been dubbed the "American Serengeti", due to the once-millions-strong herds of bison, elk, hundreds of thousands of pronghorn, as well as other now-extinct megafauna.
The present-day range of the pronghorn is generally west of the Mississippi, extending from southern Saskatchewan and Alberta, Canada south into the western US, primarily in the states of Arizona, Colorado, Idaho, Kansas, Montana, Nebraska, Nevada, New Mexico, North Dakota, Oklahoma, Oregon, South Dakota, Texas, Utah, Washington and Wyoming. In extreme Northern California, pronghorn can be found in inland counties, ranging from neighboring Nevada and Oregon, as well as the central coastal grasslands, further south. In Mexico, the Sonoran pronghorn (A. a. sonoriensis) subspecies may be found from the state of Baja California Sur east through Sonora to San Luis Potosí, in north-central regions of the country, albeit in gradually diminishing populations. They have been extirpated from Iowa and Minnesota in the United States, and from Manitoba in Canada.
Other regional subspecies include the Rocky Mountain pronghorn (A. a. americana), Mexican pronghorn (A. a. mexicana), the Oregon pronghorn (A. a. oregona), and the critically endangered Baja California pronghorn (A. a. peninsularis).
Pronghorns prefer open, expansive terrain at elevations varying between , with the densest populations in areas receiving around of rainfall per year. They eat a wide variety of plant foods, often including plants unpalatable or toxic to domestic livestock, though they also compete with them for food. In one study, forbs comprised 62% of their diet, shrubs 23%, and grasses 15%, while in another, cacti comprised 40%, grass 22%, forbs 20%, and shrubs 18%. Pronghorns also chew and eat (ruminate) cud.
Healthy pronghorn populations tend to stay within of a water source. The majority are found within of a water source.
An ongoing study by the Lava Lake Institute for Science and Conservation and the Wildlife Conservation Society shows an overland migration route that covers more than . The migrating pronghorn start travel from the foothills of the Pioneer Mountains through Craters of the Moon National Monument to the Continental Divide. Dr. Scott Bergen of the Wildlife Conservation Society says "This study shows that pronghorn are the true marathoners of the American West. With these new findings, we can confirm that Idaho supports a major overland mammal migration - an increasingly rare phenomenon in the U.S. and worldwide."
Cougars (Puma concolor), wolves (Canis lupus), coyotes (Canis latrans), grizzly bears (Ursus arctos horribilis) and bobcats (Lynx rufus) are major predators of pronghorns. Golden eagles (Aquila chrysaetos) have been reported to prey on fawns and adults. Jaguars (Panthera onca) also likely prey on pronghorns in their native range in the southwestern United States and in northern Mexico. In the Pleistocene, jaguars would likely be dangerous to pronghorns as a short-range ambush predator.
Social behavior and reproduction
Pronghorns form mixed-sex herds in the winter. In early spring, the herds break up, with young males forming bachelor groups, females forming harems, and adult males living solitarily. Some female bands share the same summer range, and bachelor male bands form between spring and fall. Females form dominance hierarchies with few circular relationships. Dominant females aggressively displace other females from feeding sites.
Adult males either defend a fixed territory that females may enter, or defend a harem of females. A pronghorn may change mating strategies depending on environmental or demographic conditions. Where precipitation is high, adult males tend to be territorial and maintain their territories with scent marking, vocalizing, and challenging intruders. In these systems, territorial males have access to better resources than bachelor males. Females also employ different mating strategies. "Sampling" females visit several males and remain with each for a short time before switching to the next male at an increasing rate as estrous approaches. "Inciting" females behave as samplers until estrous, and then incite conflicts between males, watching and then mating with the winners. Before fighting, males try to intimidate each other. If intimidation fails, they lock horns and try to injure each other. "Quiet" females remain with a single male in an isolated area throughout estrous. Females continue this mating behavior for two to three weeks.
When courting an estrous female, a male pronghorn approaches her while softly vocalizing and waving his head side to side, displaying his cheek patches. The scent glands on the pronghorn are on either side of the jaw, between the hooves, and on the rump. A receptive female remains motionless, sniffs his scent gland, and then allows the male to mount her.
Pronghorns have a gestation period of 7–8 months, which is longer than is typical for North American ungulates. They breed in mid-September, and the doe carries her fawn until late May. The gestation period is around six weeks longer than that of the white-tailed deer. Females usually bear within a few days of each other. Twin fawns are common. Newborn pronghorns weigh , most commonly . In their first 21–26 days, fawns spend time hiding in vegetation. Fawns interact with their mothers for 20–25 minutes a day; this continues even when the fawn joins a nursery. The females nurse, groom, and lead their young to food and water, as well as keep predators away from them. Females usually nurse the young about three times a day. Males are weaned 2–3 weeks earlier than females. Sexual maturity is reached at 15 to 16 months, though males rarely breed until three years old. Their lifespan is typically up to 10 years, rarely 15 years.
Relationship with humans
In regions inhabited by the Plains Indians tribes, as well as the Northwest Plateau, pronghorn was hunted as a principal food source by the local people. The pronghorn has also featured prominently in Native American mythology and oral history.
Merriwether Lewis and William Clark made several other observations on the behavior of the pronghorn and how the local tribes hunted them. They described the animal, which they referred to as the "Antelope" or the "Goat", as follows:
Population and conservation
At the turn of the 20th century, members of the wildlife conservation group Boone and Crockett Club had determined that the extinction of the pronghorn was likely. In a letter from George Bird Grinnell, Boone and Crockett Club chairman of the game preservation committee, to Walter L. Fisher, Secretary of the Interior, Grinnell stated, "The Club is much concerned about the fate of the pronghorn which appears to be everywhere rapidly diminishing." By the 1920s, hunting pressure had reduced the pronghorn population to about 13,000. Boone and Crockett Club member Charles Alexander Sheldon, in a letter to fellow member Grinnell, wrote, "Personally, I think that the antelope are doomed, yet every attempt should be made to save them." Although the club had begun their efforts to save the pronghorn in 1910 by funding and restocking the Wichita Game Refuge in Oklahoma, the National Bison Range in Montana, and the Wind Cave National Park, in South Dakota, most of the efforts were doomed since experience demonstrated that after initial increases the pronghorns would die off because of the fenced enclosures.
In 1927, Grinnell spearheaded efforts along with the help of T. Gilbert Pearson of Grinnell's National Audubon Society to create the Charles Alexander Sheldon Antelope Refuge in northern Nevada. About 2900 acres of land were jointly purchased by the two organizations and subsequently turned over to the Biological Survey as a pronghorn refuge. This donation was contingent upon the government's adding 30,000 acres of surrounding public lands. On June 20, 1929, United States President Herbert Hoover included the required public lands upon request of the Department of Agriculture and the Department of the Interior after learning that the Boone and Crockett Club and the National Audubon Society were underwriting the private land buyout. On January 26, 1931, Hoover signed the executive order for the refuge. On December 31, 1936, President Franklin Roosevelt signed an executive order creating a tract; this was the true beginning for pronghorn recovery in North America.
The protection of habitat and hunting restrictions have allowed pronghorn numbers to recover to an estimated population between 500,000 and 1,000,000 since the 1930s. Some recent decline has occurred in a few localized populations, due to bluetongue disease which is spread from sheep, but the overall trend has been positive.
Pronghorn migration corridors are threatened by habitat fragmentation and the blocking of traditional routes. In a migration study conducted by Lava Lake Institute for Science and Conservation and the Wildlife Conservation Society, at one point, the migration corridor bottlenecks to an area only 200 yards wide.
Pronghorns are now quite numerous, and outnumbered people in Wyoming and parts of northern Colorado until just recently. They are legally hunted in western states for purposes of population control and food. No major range-wide threats exist, although localized declines are taking place, particularly to the Sonoran pronghorn, mainly as a result of livestock grazing, the construction of roads, fences, and other barriers that prevent access to historical habitat, illegal hunting, insufficient forage and water, and lack of recruitment.
Three subspecies are considered endangered in all (A. a. sonoriensis, A. a. peninsularis), or part of their ranges (A. a. mexicana). The Sonoran pronghorn has an estimated population of fewer than 300 in the United States and 200–500 in Mexico, while there are approximately 200 Peninsula pronghorn in Baja California. Populations of the Sonoran pronghorn in Arizona and Mexico are protected under the Endangered Species Act (since 1967), and a recovery plan for this subspecies has been prepared by U.S. Fish and Wildlife Service. Mexican animals are listed on CITES Appendix I. Pronghorns have game-animal status in all of the western states of the United States, and permits are required to trap or hunt pronghorns.
Explanatory notes
| Biology and health sciences | Artiodactyla | null |
18952693 | https://en.wikipedia.org/wiki/Plant%20stem | Plant stem | A stem is one of two main structural axes of a vascular plant, the other being the root. It supports leaves, flowers and fruits, transports water and dissolved substances between the roots and the shoots in the xylem and phloem, engages in photosynthesis, stores nutrients, and produces new living tissue. The stem can also be called the culm, halm, haulm, stalk, or thyrsus.
The stem is normally divided into nodes and internodes:
The nodes are the points of attachment for leaves and can hold one or more leaves. There are sometimes axillary buds between the stem and leaf which can grow into branches (with leaves, conifer cones, or flowers). Adventitious roots (e.g. brace roots) may also be produced from the nodes. Vines may produce tendrils from nodes.
The internodes distance one node from another.
The term "shoots" is often confused with "stems"; "shoots" generally refers to new fresh plant growth, including both stems and other structures like leaves or flowers.
In most plants, stems are located above the soil surface, but some plants have underground stems.
Stems have several main functions:
Support for and the elevation of leaves, flowers, and fruits. The stems keep the leaves in the light and provide a place for the plant to keep its flowers and fruits.
Transport of fluids between the roots and the shoots in the xylem and phloem.
Storage of nutrients.
Production of new living tissue. The normal lifespan of plant cells is one to three years. Stems have cells called meristems that annually generate new living tissue.
Photosynthesis.
Stems have two pipe-like tissues called xylem and phloem. The xylem tissue arises from the cell facing inside and transports water by the action of transpiration pull, capillary action, and root pressure. The phloem tissue arises from the cell facing outside and consists of sieve tubes and their companion cells. The function of phloem tissue is to distribute food from photosynthetic tissue to other tissues. The two tissues are separated by cambium, a tissue that divides to form xylem or phloem cells.
Specialized terms
Stems are often specialized for storage, asexual reproduction, protection, or photosynthesis, including the following:
Acaulescent: Used to describe stems in plants that appear to be stemless. Actually these stems are just extremely short, the leaves appearing to rise directly out of the ground, e.g. some Viola species.
Arborescent: Tree with woody stems normally with a single trunk.
Axillary bud: A bud which grows at the point of attachment of an older leaf with the stem. It potentially gives rise to a shoot.
Branched: Aerial stems are described as being branched or unbranched.
Bud: An embryonic shoot with immature stem tip.
Bulb: A short vertical underground stem with fleshy storage leaves attached, e.g. onion, daffodil, and tulip. Bulbs often function in reproduction by splitting to form new bulbs or producing small new bulbs termed bulblets. Bulbs are a combination of stem and leaves so may better be considered as leaves because the leaves make up the greater part.
Caespitose: When stems grow in a tangled mass or clump or in low growing mats.
Cladode (including phylloclade): A flattened stem that appears leaf-like and is specialized for photosynthesis, e.g. cactus pads.
Climbing: Stems that cling or wrap around other plants or structures.
Corm: A short enlarged underground storage stem, e.g. taro, crocus, gladiolus.
Decumbent: A stem that lies flat on the ground and turns upwards at the ends.
Fruticose: Stems that grow shrublike with woody like habit.
Herbaceous: Non woody stems which die at the end of the growing season.
Internode: An interval between two successive nodes. It possesses the ability to elongate, either from its base or from its extremity depending on the species.
Node: A point of attachment of a leaf or a twig on the stem in seed plants. A node is a very small growth zone.
Pedicel: Stems that serve as the stalk of an individual flower in an inflorescence or infrutescence.
Peduncle: A stem that supports an inflorescence or a solitary flower.
Prickle: A sharpened extension of the stem's outer layers, e.g. roses.
Pseudostem: A false stem made of the rolled bases of leaves, which may be tall, as in banana.
Rhizome: A horizontal underground stem that functions mainly in reproduction but also in storage, e.g. most ferns, iris.
Runner: A type of stolon, horizontally growing on top of the ground and rooting at the nodes, aids in reproduction. e.g. garden strawberry, Chlorophytum comosum.
Scape: A stem that holds flowers that comes out of the ground and has no normal leaves. Hosta, lily, iris, garlic.
Stolon: A horizontal stem that produces rooted plantlets at its nodes and ends, forming near the surface of the ground.
Thorn: A modified stem with a sharpened point.
Tuber: A swollen, underground storage stem adapted for storage and reproduction, e.g. potato.
Woody: Hard textured stems with secondary xylem.
Sapwood: A woody stem, the layer of secondary phloem that surrounds the heartwood; usually active in fluid transport
Stem structure
Stem usually consist of three tissues: dermal tissue, ground tissue, and vascular tissue.
Dermal tissue covers the outer surface of the stem and usually functions to protect the stem tissue, and control gas exchange. The predominant cells of dermal tissue are epidermal cells.
Ground tissue usually consists mainly of parenchyma, collenchyma and sclerenchyma cells, and they surround vascular tissue. Ground tissue is important in aiding metabolic activities (eg. respiration, photosynthesis, transport, storage) as well as acting as structural support and forming new meristems. Most or all ground tissue may be lost in woody stems.
Vascular tissue, consisting of xylem, phloem and cambium; provides long distance transport of water, minerals and metabolites (sugars, amino acids); whilst aiding structural support and growth. The arrangement of the vascular tissues varies widely among plant species.
Dicot stems
Dicot stems with primary growth have pith in the center, with vascular bundles forming a distinct ring visible when the stem is viewed in cross section. The outside of the stem is covered with an epidermis, which is covered by a waterproof cuticle. The epidermis also may contain stomata for gas exchange and multicellular stem hairs called trichomes. A cortex consisting of hypodermis (collenchyma cells) and endodermis (starch containing cells) is present above the pericycle and vascular bundles.
Woody dicots and many nonwoody dicots have secondary growth originating from their lateral or secondary meristems: the vascular cambium and the cork cambium or phellogen. The vascular cambium forms between the xylem and phloem in the vascular bundles and connects to form a continuous cylinder. The vascular cambium cells divide to produce secondary xylem to the inside and secondary phloem to the outside. As the stem increases in diameter due to production of secondary xylem and secondary phloem, the cortex and epidermis are eventually destroyed. Before the cortex is destroyed, a cork cambium develops there. The cork cambium divides to produce waterproof cork cells externally and sometimes phelloderm cells internally. Those three tissues form the periderm, which replaces the epidermis in function. Areas of loosely packed cells in the periderm that function in gas exchange are called lenticels.
Secondary xylem is commercially important as wood. The seasonal variation in growth from the vascular cambium is what creates yearly tree rings in temperate climates. Tree rings are the basis of dendrochronology, which dates wooden objects and associated artifacts. Dendroclimatology is the use of tree rings as a record of past climates. The aerial stem of an adult tree is called a trunk. The dead, usually darker inner wood of a large diameter trunk is termed the heartwood and is the result of tylosis. The outer, living wood is termed the sapwood.
Monocot stems
Vascular bundles are present throughout the monocot stem, although concentrated towards the outside. This differs from the dicot stem that has a ring of vascular bundles and often none in the center. The shoot apex in monocot stems is more elongated. Leaf sheathes grow up around it, protecting it. This is true to some extent of almost all monocots. Monocots rarely produce secondary growth and are therefore seldom woody, with palms and bamboo being notable exceptions. However, many monocot stems increase in diameter via anomalous secondary growth.
Gymnosperm stems
All gymnosperms are woody plants. Their stems are similar in structure to woody dicots except that most gymnosperms produce only tracheids in their xylem, not the vessels found in dicots. Gymnosperm wood also often contains resin ducts. Woody dicots are called hardwoods, e.g. oak, maple and walnut. In contrast, softwoods are gymnosperms, such as pine, spruce and fir.
Fern stems
Most ferns have rhizomes with no vertical stem. The exception is tree ferns, which have vertical stems that can grow up to about 20 metres. The stem anatomy of ferns is more complicated than that of dicots because fern stems often have one or more leaf gaps in cross section. A leaf gap is where the vascular tissue branches off to a frond. In cross section, the vascular tissue does not form a complete cylinder where a leaf gap occurs. Fern stems may have solenosteles or dictyosteles or variations of them. Many fern stems have phloem tissue on both sides of the xylem in cross-section.
Relation to xenobiotics
Foreign chemicals such as air pollutants, herbicides and pesticides can damage stem structures.
Economic importance
There are thousands of species whose stems have economic uses. Stems provide a few major staple crops such as potato and taro. Sugarcane stems are a major source of sugar. Maple sugar is obtained from trunks of maple trees. Vegetables from stems are asparagus, bamboo shoots, cactus pads or nopalitos, kohlrabi, and water chestnut. The spice, cinnamon is bark from a tree trunk. Gum arabic is an important food additive obtained from the trunks of Acacia senegal trees. Chicle, the main ingredient in chewing gum, is obtained from trunks of the chicle tree.
Medicines obtained from stems include quinine from the bark of cinchona trees, camphor distilled from wood of a tree in the same genus that provides cinnamon, and the muscle relaxant curare from the bark of tropical vines.
Wood is used in thousands of ways; it can be used to create buildings, furniture, boats, airplanes, wagons, car parts, musical instruments, sports equipment, railroad ties, utility poles, fence posts, pilings, toothpicks, matches, plywood, coffins, shingles, barrel staves, toys, tool handles, picture frames, veneer, charcoal and firewood. Wood pulp is widely used to make paper, paperboard, cellulose sponges, cellophane and some important plastics and textiles, such as cellulose acetate and rayon. Bamboo stems also have hundreds of uses, including in paper, buildings, furniture, boats, musical instruments, fishing poles, water pipes, plant stakes, and scaffolding. Trunks of palms and tree ferns are often used for building. Stems of reed are an important building material for use in thatching in some areas.
Tannins used for tanning leather are obtained from the wood of certain trees, such as quebracho. Cork is obtained from the bark of the cork oak. Rubber is obtained from the trunks of Hevea brasiliensis. Rattan, used for furniture and baskets, is made from the stems of tropical vining palms. Bast fibers for textiles and rope are obtained from stems of plants like flax, hemp, jute and ramie. The earliest known paper was obtained from the stems of papyrus by the ancient Egyptians.
Amber is fossilized sap from tree trunks; it is used for jewelry and may contain preserved animals. Resins from conifer wood are used to produce turpentine and rosin. Tree bark is often used as a mulch and in growing media for container plants. It also can become the natural habitat of lichens.
Some ornamental plants are grown mainly for their attractive stems, e.g.:
White bark of paper birch
Twisted branches of corkscrew willow and Harry Lauder's walking stick (Corylus avellana 'Contorta')
Red, peeling bark of paperbark maple
| Biology and health sciences | Plant stem | null |
18952739 | https://en.wikipedia.org/wiki/Carnivorous%20plant | Carnivorous plant | Carnivorous plants are plants that derive some or most of their nutrients from trapping and consuming animals or protozoans, typically insects and other arthropods, and occasionally small mammals and birds. They still generate all of their energy from photosynthesis. They have adapted to grow in waterlogged sunny places where the soil is thin or poor in nutrients, especially nitrogen, such as acidic bogs.
They can be found on all continents except Antarctica, as well as many Pacific islands. In 1875, Charles Darwin published Insectivorous Plants, the first treatise to recognize the significance of carnivory in plants, describing years of painstaking research.
True carnivory is believed to have evolved independently at least 12 times in five different orders of flowering plants, and is represented by more than a dozen genera. This classification includes at least 583 species that attract, trap, and kill prey, absorbing the resulting available nutrients. Venus flytraps (Dionaea muscipula), pitcher plants, and bladderworts (Utricularia spp.) can be seen as exemplars of key traits genetically associated with carnivory: trap leaf development, prey digestion, and nutrient absorption.
The number of known species has increased by approximately 3 species per year since the year 2000. Additionally, over 300 protocarnivorous plant species in several genera show some but not all of these characteristics. A 2020 assessment has found that roughly one quarter are threatened with extinction from human actions.
Definition
Plants are considered carnivorous if they have these five traits:
capture prey in traps
kill the captured prey
digest the captured prey
absorb nutrients from the killed and digested prey
use those nutrients to grow and develop.
Other traits may include the attraction and retention of prey.
Trapping mechanisms
Five basic trapping mechanisms are found in carnivorous plants.
Pitfall traps (pitcher plants) trap prey in a rolled leaf that contains a pool of digestive enzymes or bacteria.
Flypaper traps use a sticky mucilage.
Snap traps utilise rapid leaf movements.
Bladder traps suck in prey with a bladder that generates an internal vacuum.
Lobster-pot traps, also known as eel traps, use inward-pointing hairs to force prey to move towards a digestive organ.
These traps may be active or passive, depending on whether movement aids the capture of prey. For example, Triphyophyllum is a passive flypaper that secretes mucilage, but whose leaves do not grow or move in response to prey capture. Meanwhile, sundews are active flypaper traps whose leaves undergo rapid acid growth, which is an expansion of individual cells as opposed to cell division. The rapid acid growth allows the sundews' tentacles to bend, aiding in the retention and digestion of prey.
Pitfall traps
Characterised by an internal chamber, pitfall traps are thought to have evolved independently at least six times. This particular adaptation is found within the families Sarraceniaceae (Darlingtonia, Heliamphora, Sarracenia), Nepenthaceae (Nepenthes), and Cephalotaceae (Cephalotus). Within the family Bromeliaceae, pitcher morphology and carnivory evolved twice (Brocchinia and Catopsis). Because these families do not share a common ancestor who also had pitfall trap morphology, carnivorous pitchers are an example of convergent evolution.
A passive trap, pitfall traps attract prey with nectar bribes secreted by the peristome and bright flower-like anthocyanin patterning within the pitcher. The linings of most pitcher plants are covered in a loose coating of waxy flakes which are slippery for insects, causing them to fall into the pitcher. Once within the pitcher structure, digestive enzymes or mutualistic species break down the prey into an absorbable form for the plant. Water can become trapped within the pitcher, making a habitat for other flora and fauna. This type of 'water body' is called a phytotelma.
The simplest pitcher plants are probably those of Heliamphora, the marsh pitcher plant. In this genus, the traps are clearly derived from a simple rolled leaf whose margins have sealed together. These plants live in areas of high rainfall in South America such as Mount Roraima and consequently have a problem ensuring their pitchers do not overflow. To counteract this problem, natural selection has favoured the evolution of an overflow similar to that of a bathroom sink—a small gap in the zipped-up leaf margins allows excess water to flow out of the pitcher.
In the genus Sarracenia, the problem of pitcher overflow is solved by an operculum, which is essentially a flared leaflet that covers the opening of the rolled-leaf tube and protects it from rain. Possibly because of this improved waterproofing, Sarracenia species secrete enzymes such as proteases and phosphatases into the digestive fluid at the bottom of the pitcher. In at least one species, Sarracenia flava, the nectar bribe is laced with coniine, a toxic alkaloid also found in hemlock, which probably increases the efficiency of the traps by intoxicating prey.
Most Heliamphora rely on bacterial digestion alone with the exception of a single species, Heliamphora tatei, which does produce digestive enzymes. The enzymes digest the proteins and nucleic acids in the prey, releasing amino acids and phosphate ions, which the plant absorbs.
Darlingtonia californica, the cobra plant, possesses an adaptation also found in Sarracenia psittacina and, to a lesser extent, in Sarracenia minor: the operculum is balloon-like and almost seals the opening to the tube. This balloon-like chamber is pitted with areolae, chlorophyll-free patches through which light can penetrate. Insects, mostly ants, enter the chamber via the opening underneath the balloon. Once inside, they tire themselves trying to escape from these false exits, until they eventually fall into the tube. Prey access is increased by the "fish tails", outgrowths of the operculum that give the plant its name. Some seedling Sarracenia species also have long, overhanging opercular outgrowths; Darlingtonia may therefore represent an example of neoteny.
The second major group of pitcher plants are the monkey cups or tropical pitcher plants of the genus Nepenthes. In the hundred or so species of this genus, the pitcher is borne at the end of a tendril, which grows as an extension to the midrib of the leaf. Most species catch insects, although the larger ones, such as Nepenthes rajah, also occasionally take small mammals and reptiles. Nepenthes bicalcarata possesses two sharp thorns that project from the base of the operculum over the entrance to the pitcher. These likely serve to lure insects into a precarious position over the pitcher mouth, where they may lose their footing and fall into the fluid within.
The pitfall trap has evolved independently in at least two other groups. The Albany pitcher plant, Cephalotus follicularis, is a small pitcher plant from Western Australia, with moccasin-like pitchers. The rim of its pitcher's opening (the peristome) is particularly pronounced (both secrete nectar) and provides a thorny overhang to the opening, preventing trapped insects from climbing out.
The final carnivore with a pitfall-like trap is the bromeliad Brocchinia reducta. Like most relatives of the pineapple, the tightly packed, waxy leaf bases of the strap-like leaves of this species form an urn. In most bromeliads, water collects readily in this urn and may provide habitats for frogs, insects and, more useful for the plant, diazotrophic (nitrogen-fixing) bacteria. In Brocchinia, the urn is a specialised insect trap, with a loose, waxy lining and a population of digestive bacteria.
Flypaper traps
The flypaper trap utilises sticky mucilage or glue. The leaf of flypaper traps is studded with mucilage-secreting glands, which may be short (like those of the butterworts), or long and mobile (like those of many sundews). Flypapers have evolved independently at least five times. There is evidence that some clades of flypaper traps have evolved from morphologically more complex traps such as pitchers.
In the genus Pinguicula, the mucilage glands are quite short (sessile), and the leaf, while shiny (giving the genus its common name of 'butterwort'), does not appear carnivorous. However, this belies the fact that the leaf is an extremely effective trap of small flying insects (such as fungus gnats), and its surface responds to prey by relatively rapid growth. This thigmotropic growth may involve rolling of the leaf blade (to prevent rain from splashing the prey off the leaf surface) or dishing of the surface under the prey to form a shallow digestive pit.
The sundew genus (Drosera) consists of over 100 species of active flypapers whose mucilage glands are borne at the end of long tentacles, which frequently grow fast enough in response to prey (thigmotropism) to aid the trapping process. The tentacles of D. burmanii can bend 180° in a minute or so. Sundews are extremely cosmopolitan and are found on all the continents except the Antarctic mainland. They are most diverse in Australia, the home to the large subgroup of pygmy sundews such as D. pygmaea and to a number of tuberous sundews such as D. peltata, which form tubers that aestivate during the dry summer months. These species are so dependent on insect sources of nitrogen that they generally lack the enzyme nitrate reductase, which most plants require to assimilate soil-borne nitrate into organic forms.
Similar to Drosera is the Portuguese dewy pine, Drosophyllum, which differs from the sundews in being passive. Its leaves are incapable of rapid movement or growth. Unrelated, but similar in habit, are the Australian rainbow plants (Byblis). Drosophyllum is unusual in that it grows under near-desert conditions; almost all other carnivores are either bog plants or grow in moist tropical areas.
Recent molecular data (particularly the production of plumbagin) indicate that the remaining flypaper, Triphyophyllum peltatum, a member of the Dioncophyllaceae, is closely related to Drosophyllum and forms part of a larger clade of carnivorous and non-carnivorous plants with the Droseraceae, Nepenthaceae, Ancistrocladaceae and Plumbaginaceae. This plant is usually encountered as a liana, but in its juvenile phase, the plant is carnivorous. This may be related to a requirement for specific nutrients for flowering.
Snap traps
The only two active snap traps—the Venus flytrap (Dionaea muscipula) and the waterwheel plant (Aldrovanda vesiculosa)—had a common ancestor with the snap trap adaptation, which had evolved from an ancestral lineage that utilized flypaper traps. Their trapping mechanism has also been described as a "mouse trap", "bear trap" or "man trap", based on their shape and rapid movement. However, the term snap trap is preferred as other designations are misleading, particularly with respect to the intended prey. Aldrovanda is aquatic and specialized in catching small invertebrates; Dionaea is terrestrial and catches a variety of arthropods, including spiders.
The traps are very similar, with leaves whose terminal section is divided into two lobes, hinged along the midrib. Trigger hairs (three on each lobe in Dionaea muscipula, many more in the case of Aldrovanda) inside the trap lobes are sensitive to touch. When a trigger hair is bent, stretch-gated ion channels in the membranes of cells at the base of the trigger hair open, generating an action potential that propagates to cells in the midrib. These cells respond by pumping out ions, which may either cause water to follow by osmosis (collapsing the cells in the midrib) or cause rapid acid growth. The mechanism is still debated, but in any case, changes in the shape of cells in the midrib allow the lobes, held under tension, to snap shut, flipping rapidly from convex to concave and interring the prey. This whole process takes less than a second. In the Venus flytrap, closure in response to raindrops and blown-in debris is prevented by the leaves having a simple memory: for the lobes to shut, two stimuli are required, 0.5 to 30 seconds apart. According to a recent study, calcium molecules move dynamically within the cells of the plant's leaves when a carnivorous plant touches live prey. Changing calcium levels make leaves move to catch prey, likely by producing more hormones related to defense.
The snapping of the leaves is a case of thigmonasty (undirected movement in response to touch). Further stimulation of the lobe's internal surfaces by the struggling insects causes the lobes to close even tighter (thigmotropism), sealing the lobes hermetically and forming a stomach in which digestion occurs over a period of one to two weeks. Once this process is triggered, it cannot be reversed and requires more stimulation to trigger the next steps. Leaves can be reused three or four times before they become unresponsive to stimulation, depending on the growing conditions.
Bladder traps
Bladder traps are exclusive to the genus Utricularia, or bladderworts. The bladders (vesiculae) pump ions out of their interiors. Water follows by osmosis, generating a partial vacuum inside the bladder. The bladder has a small opening, sealed by a hinged door. In aquatic species, the door has a pair of long trigger hairs. Aquatic invertebrates such as Daphnia touch these hairs and deform the door by lever action, releasing the vacuum. The invertebrate is sucked into the bladder, where it is digested. Many species of Utricularia (such as U. sandersonii) are terrestrial, growing in waterlogged soil, and their trapping mechanism is triggered in a slightly different manner. Bladderworts lack roots, but terrestrial species have anchoring stems that resemble roots. Temperate aquatic bladderworts generally die back to a resting turion during the winter months, and U. macrorhiza appears to regulate the number of bladders it bears in response to the prevailing nutrient content of its habitat.
Lobster-pot traps
A lobster-pot trap is a chamber that is easy to enter, and whose exit is either difficult to find or obstructed by inward-pointing bristles. Lobster pots are the trapping mechanism in Genlisea, the corkscrew plants. These plants appear to specialise in aquatic protozoa. A Y-shaped modified leaf allows prey to enter but not exit. Inward-pointing hairs force the prey to move in a particular direction. Prey entering the spiral entrance that coils around the upper two arms of the Y are forced to move inexorably towards a stomach in the lower arm of the Y, where they are digested. Prey movement is also thought to be encouraged by water movement through the trap, produced in a similar way to the vacuum in bladder traps, and probably evolutionarily related to it.
Outside of Genlisea, features reminiscent of lobster-pot traps can be seen in Sarracenia psittacina, Darlingtonia californica, and, some horticulturalists argue, Nepenthes aristolochioides.
Combination traps
The trapping mechanism of the sundew Drosera glanduligera combines features of both flypaper and snap traps; it has been termed a catapult-flypaper trap. Similarly, Nepenthes jamban is a combination of pitfall and flypaper traps because it has a sticky pitcher fluid.
Most Sumatran nepenthes, like N. inermis, also have this method. For example, N. dubia and N. flava also use this method.
Borderline carnivores
To be defined as carnivorous, a plant must first exhibit an adaptation of some trait specifically for the attraction, capture, or digestion of prey. Only one trait needs to have evolved that fits this adaptive requirement, as many current carnivorous plant genera lack some of the above-mentioned attributes. The second requirement is the ability to absorb nutrients from dead prey and gain a fitness advantage from the integration of these derived nutrients (mostly amino acids and ammonium ions) either through increased growth or pollen and/or seed production. However, plants that may opportunistically utilise nutrients from dead animals without specifically seeking and capturing fauna are excluded from the carnivorous definition. The second requirement also differentiates carnivory from defensive plant characteristics that may kill or incapacitate insects without the advantage of nutrient absorption. Due to the observation that many currently classified carnivores lack digestive enzymes for breaking down nutrients and instead rely upon mutualistic and symbiotic relationships with bacteria, ants, or insects, this adaptation has been added to the carnivorous definition. Despite this, there are cases where plants appear carnivorous, in that they fulfill some of the above definition, but are not truly carnivorous. Some botanists argue that there is a spectrum of carnivory found in plants: from completely non-carnivorous plants like cabbages, to borderline carnivores, to unspecialised and simple traps, like Heliamphora, to extremely specialised and complex traps, like that of the Venus flytrap.
A possible carnivore is the genus Roridula; the plants in this genus produce sticky leaves with resin-tipped glands and look extremely similar to some of the larger sundews. However, they do not directly benefit from the insects they catch. Instead, they form a mutualistic symbiosis with species of assassin bug (genus Pameridea), which eat the trapped insects. The plant benefits from the nutrients in the bugs' feces. By some definitions this would still constitute botanical carnivory.
A number of species in the Martyniaceae (previously Pedaliaceae), such as Ibicella lutea, have sticky leaves that trap insects. However, these plants have not been shown conclusively to be carnivorous. Likewise, the seeds of Shepherd's Purse, urns of Paepalanthus bromelioides, bracts of Passiflora foetida, and flower stalks and sepals of triggerplants (Stylidium) appear to trap and kill insects, but their classification as carnivores is contentious.
Two genera of liverwort, Colura and Pleurozia, have sac-shaped leaves that trap and kill ciliates and may digest them.
A species of pitcher plant, Nepenthes ampullaria, has evolved away from being a carnivore. Rather than catching animals, it catches falling leaves in its pitchers.
Digestion
Specialized multicellular secretion glands produce digestive fluid that smother, kill, and digest prey as well as make a solution to assimilate released nutrients. Saccharides are often found in plants that have adhesive traps or plants that use viscous secretion to retain captured prey. The digestion fluid is often nutrient poor and has ions K, Na, Ca and Mg (for species in the Nepenthes genera for example), along with numerous proteins which vary across genera. Peroxidases are also involved for some species. The body of the prey is decomposed by a cocktail of hydrolytic enzymes which are stored in sub-cellular compartments or synthesized over and over as needed.
Proteins of digestive fluid include proteases, chitinases (partly destroy exoskeleton of insects), phosphatases, and nucleases.
Evolution
General pattern of independent development in multiple lineages
Charles Darwin spent 16 years growing carnivorous plants, experimenting with them in the greenhouse of his home in Kent, Down House. In his pioneering book Insectivorous Plants (1875) Darwin concluded that carnivory in plants was convergent, writing that carnivorous genera Utricularia and Nepenthes were not "at all related to the [carnivorous family] Droseraceae". This remained a subject of debate for over a century. In 1960, Leon Croizat concluded that carnivory was monophyletic, and placed all the carnivorous plants together at the base of the angiosperms. Molecular studies over the past 30 years have led to a wide consensus that Darwin was correct, with studies showing that carnivory evolved at least six times in the angiosperms, and that trap designs such as pitcher traps and flypaper traps are analogous rather than homologous.
Researchers using molecular data have concluded that carnivory evolved independently in the Poales (Brocchinia and Catopsis in the Bromeliaceae), the Caryophyllales (Droseraceae, Nepenthaceae, Drosophyllaceae, Dioncophyllaceae), the Oxalidales (Cephalotus), the Ericales (Sarraceniaceae and Roridulaceae), and twice in the Lamiales (Lentibulariaceae and independently in Byblidaceae). The oldest evolution of an existing carnivory lineage has been dated to 85.6 million years ago, with the most recent being Brocchinia reducta in the Bromeliaceae estimated at only 1.9 mya.
The evolution of carnivorous plants is obscured by the paucity of their fossil record. Very few fossils have been found, and then usually only as seed or pollen. Carnivorous plants are generally herbs, and their traps are produced by primary growth. They generally do not form readily fossilisable structures such as thick bark or wood.
Researchers are increasingly using genome sequencing technology to examine the development of carnivorous species and relationships between them. Genetic evidence suggests that carnivory developed by co-opting and repurposing existing genes which had established functions in flowering plants, rather than by "hijacking" genes from other types of organisms.
Adaption to extreme habitats
Most carnivorous plants live in habitats with high light, waterlogged soils, and extremely low soil nitrogen and phosphorus, producing the ecological impetus to derive nitrogen from an alternate source. High-light environments allowed for the trade-off between photosynthetic leaves and photosynthetically inefficient, prey-capturing traps. To compensate for the photosynthetically inefficient material, the nutrients obtained through carnivory would need to increase photosynthesis by investing in more leaf mass (i.e. growth). Consequently, when there is a shortage of nutrients, sufficient light and water, the capture and digestion of prey has the greatest impact on photosynthetic gains, thus favoring the evolution of plant adaptations which allowed for more effective, efficient carnivory.
Due to the required energy and resource allocations for carnivorous adaptations (e.g. the production of lures, digestive enzymes, modified leaf structures, and the decreased rate of photosynthesis over total leaf area), some authors argue that carnivory is an evolutionary "last resort" when nitrogen and phosphorus are extremely limited in an ecosystem.
Inferences from trap mechanism
Despite meager fossil evidence, much can be deduced from the structure of current traps and their ecological interactions. It is widely believed that carnivory evolved under extremely nutrient-poor conditions, leading to a cost-benefit model for botanical carnivory. Cost-benefit models are used under the assumption that there is a set amount of potential energy available to an organism, which leads to trade-offs wherein energy is allocated to certain functions to maximize competitive ability and fitness. For carnivory, the trait could only evolve if the increase in nutrients from capturing prey exceeded the cost of investment in carnivorous adaptations.
Pitfall traps are derived from rolled leaves, which evolved several independent times through convergent evolution. The vascular tissues of Sarracenia is a case in point. The keel along the front of the trap contains a mixture of leftward- and rightward-facing vascular bundles, as would be predicted from the fusion of the edges of an adaxial (stem-facing) leaf surface. Flypapers also show a simple evolutionary gradient from sticky, non-carnivorous leaves, through passive flypapers to active forms. Molecular data show the Dionaea–Aldrovanda clade is closely related to Drosera, and evolved from active flypaper traps into snap traps.
Hypothetical common start with a sticky, hairy leaf
It has been suggested that all trap types are modifications of a similar basic structure: the hairy leaf. Hairy (or more specifically, stalked-glandular) leaves can catch and retain drops of rainwater, especially if shield-shaped or peltate, thus promoting bacteria growth. Insects land on the leaf, become mired by the surface tension of the water, and suffocate. Bacteria jumpstart decay, releasing from the corpse nutrients that the plant can absorb through its leaves. This foliar feeding can be observed in most non-carnivorous plants.
Plants that were better at retaining insects or water therefore had a selective advantage. Rainwater can be retained by cupping the leaf, and pitfall traps may have evolved simply by selection pressure for the production of more deeply cupped leaves, followed by "zipping up" of the margins and subsequent loss of most of the hairs. Alternatively, insects can be retained by making the leaf stickier by the production of mucilage, leading to flypaper traps.
The only traps that are unlikely to have descended from a hairy leaf or sepal are the carnivorous bromeliads (Brocchinia and Catopsis): These plants use the urn – a characteristic part of all bromeliads, not just the carnivorous ones – for a new purpose, and build on it by the production of wax and the other paraphernalia of carnivory.
Leaves shaped like pitchers and lobster-pots
The lobster-pot traps of Genlisea are difficult to interpret. They may have developed from bifurcated pitchers that later specialised on ground-dwelling prey; or, perhaps, the prey-guiding protrusions of bladder traps became more substantial than the net-like funnel found in most aquatic bladderworts. Whatever their origin, the helical shape of the lobster pot is an adaptation that displays as much trapping surface as possible in all directions when buried in moss.
The traps of the bladderworts may have derived from pitchers that specialised in aquatic prey when flooded, like Sarracenia psittacina does today. Escaping prey in terrestrial pitchers have to climb or fly out of a trap, and both of these can be prevented by wax, gravity and narrow tubes. However, a flooded trap can be swum out of, so in Utricularia, a one-way lid may have developed to form the door of a proto-bladder. Later, this may have become active by the evolution of a partial vacuum inside the bladder, tripped by prey brushing against trigger hairs on the door of the bladder.
The active glue traps use rapid plant movements to trap their prey. Rapid plant movement can result from actual growth, or from rapid changes in cell turgor, which allow cells to expand or contract by quickly altering their water content. Slow-moving flypapers like Pinguicula exploit growth, while the Venus flytrap uses such rapid turgor changes which make glue unnecessary. The stalked glands that once made glue became teeth and trigger hairs in species with active snap traps – an example of natural selection hijacking preexisting structures for new functions.
Unclear clustering of carnivory in Caryophyllales
Recent taxonomic analysis of the relationships within the Caryophyllales indicate that the Droseraceae, Triphyophyllum, Nepenthaceae and Drosophyllum, while closely related, are embedded within a larger clade that includes non-carnivorous groups such as the tamarisks, Ancistrocladaceae, Polygonaceae and Plumbaginaceae.
The tamarisks possess specialised salt-excreting glands on their leaves, as do several of the Plumbaginaceae (such as the sea lavender, Limonium), which may have been co-opted for the excretion of other chemicals, such as proteases and mucilage. Some of the Plumbaginaceae (e.g. Ceratostigma) also have stalked, vascularised glands that secrete mucilage on their calyces and aid in seed dispersal and possibly in protecting the flowers from crawling parasitic insects. The balsams (such as Impatiens), which are closely related to the Sarraceniaceae and Roridula, similarly possess stalked glands.
Philcoxia is unique in the Plantaginaceae as a result of its subterranean stems and leaves, which have been shown to be used in the capture of nematodes. These plants grow in sand in Brazil, where they are likely to receive other nutrients. Like many other types of carnivorous plant, stalked glands are seen on the leaves. Enzymes on the leaves are used to digest the worms and release their nutrients.
Carnivory in angiosperms
Botanical carnivory has evolved in several independent families peppered throughout the angiosperm phylogeny, showing that carnivorous traits underwent convergent evolution multiple times to create similar morphologies across disparate families. Results of genetic testing published in 2017 found an example of convergent evolution – a digestive enzyme with the same functional mutations across unrelated lineages.
Ecology and modeling of carnivory
Carnivorous plants are widespread but rather rare. They are almost entirely restricted to habitats such as bogs, where soil nutrients are extremely limiting, but where sunlight and water are readily available. Only under such extreme conditions is carnivory favored to an extent that makes the adaptations advantageous.
The archetypal carnivore, the Venus flytrap, grows in soils with almost immeasurable nitrate and calcium levels. Plants need nitrogen for protein synthesis, calcium for cell wall stiffening, phosphate for nucleic acid synthesis, and iron and magnesium for chlorophyll synthesis. The soil is often waterlogged, which favours the production of toxic ions such as ammonium, and its pH is an acidic 4 to 5. Ammonium can be used as a source of nitrogen by plants, but its high toxicity means that concentrations high enough to fertilise are also high enough to cause damage.
However, the habitat is warm, sunny, constantly moist, and the plant experiences relatively little competition from low growing Sphagnum moss. Still, carnivores are also found in very atypical habitats. Drosophyllum lusitanicum is found around desert edges and Pinguicula valisneriifolia on limestone (calcium-rich) cliffs.
In all the studied cases, carnivory allows plants to grow and reproduce using animals as a source of nitrogen, phosphorus and possibly potassium. However, there is a spectrum of dependency on animal prey. Pygmy sundews are unable to use nitrate from soil because they lack the necessary enzymes (nitrate reductase in particular). Common butterworts (Pinguicula vulgaris) can use inorganic sources of nitrogen better than organic sources, but a mixture of both is preferred. European bladderworts seem to use both sources equally well. Animal prey makes up for differing deficiencies in soil nutrients.
Plants use their leaves to intercept sunlight. The energy is used to reduce carbon dioxide from the air with electrons from water to make sugars (and other biomass) and a waste product, oxygen, in the process of photosynthesis. Leaves also respire, in a similar way to animals, by burning their biomass to generate chemical energy. This energy is temporarily stored in the form of ATP (adenosine triphosphate), which acts as an energy currency for metabolism in all living things. As a waste product, respiration produces carbon dioxide.
For a plant to grow, it must photosynthesise more than it respires. Otherwise, it will eventually exhaust its biomass and die. The potential for plant growth is net photosynthesis, the total gross gain of biomass by photosynthesis, minus the biomass lost by respiration. Understanding carnivory requires a cost-benefit analysis of these factors.
In carnivorous plants, the leaf is not just used to photosynthesise, but also as a trap. Changing the leaf shape to make it a better trap generally makes it less efficient at photosynthesis. For example, pitchers have to be held upright, so that only their opercula directly intercept light. The plant also has to expend extra energy on non-photosynthetic structures like glands, hairs, glue and digestive enzymes. To produce such structures, the plant requires ATP and respires more of its biomass. Hence, a carnivorous plant will have both decreased photosynthesis and increased respiration, making the potential for growth small and the cost of carnivory high.
Being carnivorous allows the plant to grow better when the soil contains little nitrate or phosphate. In particular, an increased supply of nitrogen and phosphorus makes photosynthesis more efficient, because photosynthesis depends on the plant being able to synthesise very large amounts of the nitrogen-rich enzyme RuBisCO (ribulose-1,5-bis-phosphate carboxylase/oxygenase), the most abundant protein on Earth.
It is intuitively clear that the Venus flytrap is more carnivorous than Triphyophyllum peltatum. The former is a full-time moving snap-trap; the latter is a part-time, non-moving flypaper. The energy "wasted" by the plant in building and fuelling its trap is a suitable measure of the carnivory of the trap.
Using this measure of investment in carnivory, a model can be proposed. Above is a graph of carbon dioxide uptake (potential for growth) against trap respiration (investment in carnivory) for a leaf in a sunny habitat containing no soil nutrients at all. Respiration is a straight line sloping down under the horizontal axis (respiration produces carbon dioxide). Gross photosynthesis is a curved line above the horizontal axis: as investment increases, so too does the photosynthesis of the trap, as the leaf receives a better supply of nitrogen and phosphorus. Eventually another factor (such as light intensity or carbon dioxide concentration) will become more limiting to photosynthesis than nitrogen or phosphorus supply. As a result, increasing the investment will not make the plant grow better. The net uptake of carbon dioxide, and therefore, the plant's potential for growth, must be positive for the plant to survive. There is a broad span of investment where this is the case, and there is also a non-zero optimum. Plants investing more or less than this optimum will take up less carbon dioxide than an optimal plant, hence grow less well. These plants will be at a selective disadvantage. At zero investment the growth is zero, because a non-carnivorous plant cannot survive in a habitat with absolutely no soil-borne nutrients. Such habitats do not exist, so for example, Sphagnum absorbs the tiny amounts of nitrates and phosphates in rain very efficiently and also forms symbioses with diazotrophic cyanobacteria.
In a habitat with abundant soil nutrients but little light (as shown above), the gross photosynthesis curve will be lower and flatter, because light will be more limiting than nutrients. A plant can grow at zero investment in carnivory; this is also the optimum investment for a plant, as any investment in traps reduces net photosynthesis (growth) to less than the net photosynthesis of a plant that obtains its nutrients from soil alone.
Carnivorous plants exist between these two extremes: the less limiting light and water are, and the more limiting soil nutrients are, the higher the optimum investment in carnivory, and hence the more obvious the adaptations will be to the casual observer.
The most obvious evidence for this model is that carnivorous plants tend to grow in habitats where water and light are abundant and where competition is relatively low: the typical bog. Those that do not tend to be even more fastidious in some other way. Drosophyllum lusitanicum grows where there is little water, but it is even more extreme in its requirement for bright light and low disturbance than most other carnivores. Pinguicula valisneriifolia grows in soils with high levels of calcium but requires strong illumination and lower competition than many butterworts.
In general, carnivorous plants are poor competitors, because they invest too heavily in structures that have no selective advantage in nutrient-rich habitats. They succeed only where other plants fail. Carnivores are to nutrients what cacti are to water. Carnivory only pays off when the nutrient stress is high and where light is abundant. When these conditions are not met, some plants give up carnivory temporarily. Sarracenia spp. produce flat, non-carnivorous leaves (phyllodes) in winter. Light levels are lower than in summer, so light is more limiting than nutrients, and carnivory does not pay. The lack of insects in winter exacerbates the problem. Damage to growing pitcher leaves prevents them from forming proper pitchers, and again, the plant produces a phyllode instead.
Many other carnivores shut down in some seasons. Tuberous sundews die back to tubers in the dry season, bladderworts to turions in winter, and non-carnivorous leaves are made by most butterworts and Cephalotus in the less favourable seasons. Utricularia macrorhiza varies the number of bladders it produces based on the expected density of prey. Part-time carnivory in Triphyophyllum peltatum may be due to an unusually high need for potassium at a certain point in the life cycle, just before flowering.
The more carnivorous a plant is, the less conventional its habitat is likely to be. Venus flytraps live in a very specialised habitat, whereas less carnivorous plants (Byblis, Pinguicula) are found in less unusual habitats (i.e., those typical for non-carnivores). Byblis and Drosophyllum both come from relatively arid regions and are both passive flypapers, arguably the lowest maintenance form of trap. Venus flytraps filter their prey using the teeth around the trap's edge, so as not to waste energy on hard-to-digest prey. In evolution, laziness pays, because energy can be used for reproduction, and short-term benefits in reproduction will outweigh long-term benefits in anything else.
Carnivory rarely pays, so even carnivorous plants avoid it when there is too little light or an easier source of nutrients, and they use as few carnivorous features as are required at a given time or for a given prey item. There are very few habitats stressful enough to make investing biomass and energy in trigger hairs and enzymes worthwhile. Many plants occasionally benefit from animal protein rotting on their leaves, but carnivory that is obvious enough for the casual observer to notice is rare.
Bromeliads seem very well preadapted to carnivory, but only one or two species can be classified as truly carnivorous. By their very shape, bromeliads will benefit from increased prey-derived nutrient input. In this sense, bromeliads are probably carnivorous, but their habitats are too dark for more extreme, recognisable carnivory to evolve. Most bromeliads are epiphytes, and most epiphytes grow in partial shade on tree branches. Brocchinia reducta, on the other hand, is a ground dweller.
Many carnivorous plants are not strongly competitive and rely on circumstances to suppress dominating vegetation. Accordingly, some of them rely on fire ecology for their continued survival.
For the most part carnivorous plant populations are not dominant enough to be dramatically significant, ecologically speaking, but there is an impressive variety of organisms that interact with various carnivorous plants in sundry relationships of kleptoparasitism, commensalism, and mutualism. For example, small insectivores such as tree frogs often exploit the supply of prey to be found in pitcher plants, and the frog Microhyla nepenthicola actually specialises in such habitats. Certain crab spiders such as Henriksenia nepenthicola and H. labuanica live largely on the prey of Nepenthes, and other, less specialised, spiders may build webs where they trap insects attracted by the smell or appearance of the traps; some scavengers, detritivores, and also organisms that harvest or exploit those in turn, such as the mosquito Wyeomyia smithii are largely or totally dependent on particular carnivorous plants. Plants such as Roridula species combine with specialised bugs (Pameridea roridulae) in benefiting from insects trapped on their leaves.
Associations with species of pitcher plants are so many and varied that the study of Nepenthes infauna is something of a discipline in its own right. Camponotus schmitzi, the diving ant, has an intimate degree of mutualism with the pitcher plant Nepenthes bicalcarata; it not only retrieves prey and detritus from beneath the surface of the liquid in the pitchers, but repels herbivores, and cleans the pitcher peristome, maintaining its slippery nature. The ants have been reported to attack struggling prey, hindering their escape, so there might be an element of myrmecotrophy to the relationship. Numerous species of mosquitoes lay their eggs in the liquid, where their larvae play various roles, depending on species; some eat microbes and detritus, as is common among mosquito larvae, whereas some species of Toxorhynchites also breed in pitchers, and their larvae are predators of other species of mosquito larvae. Apart from the crab spiders on pitchers, an actual small, red crab Geosesarma malayanum will enter the fluid, robbing and scavenging, though reputedly it does so at some risk of being captured and digested itself.
Nepenthes rajah has a remarkable mutualism with two unrelated small mammals, the mountain treeshrew (Tupaia montana) and the summit rat (Rattus baluensis). The tree shrews and the rats defecate into the plant's traps while visiting them to feed on sweet, fruity secretions from glands on the pitcher lids. The tree shrew also has a similar relationship with at least two other giant species of Nepenthes. More subtly, Hardwicke's woolly bat (Kerivoula hardwickii), a small species, roosts beneath the operculum (lid) of Nepenthes hemsleyana. The bat's excretions that land in the pitcher pay for the shelter, as it were. To the plant the excreta are more readily assimilable than intact insects would be.
There also is a considerable list of Nepenthes endophytes; these are microbes other than pathogens that live in the tissues of pitcher plants, often apparently harmlessly.
Another important area of symbiosis between carnivorous plants and insects is pollination. While many species of carnivorous plant can reproduce asexually via self-pollination or vegetative propagation, many carnivorous plants are insect-pollinated. Outcross pollination is beneficial as it increases genetic diversity. This means that carnivorous plants undergo an evolutionary and ecological conflict often called the pollinator-prey conflict. There are several ways by which carnivorous plants reduce the strain of the pollinator-prey conflict. For long-lived plants, the short-term loss of reproduction may be offset by the future growth made possible by resources obtained from prey. Other plants might "target" different species of insect for pollination and prey using different olfactory and visual cues.
Conservation threats
Approximately half of the plant species assessed by the IUCN are considered threatened (vulnerable, endangered or critically endangered). Common threats are habitat loss as a result of agriculture, collection of wild plants, pollution, invasive species, residential and commercial development, energy production, mining, transportation services, geologic events, climate change, severe weather, and many other anthropogenic activities. Species in the same genus were proven to face similar threats. Threat by continent is deemed highly variable, with threats found for 19 species in North America, 15 species in Asia, seven species in Europe, six species in South America, two species in Africa, and one species in Australia Indicator species' such as Sarracenia reveal positive associations with regard to these threats. Certain threats are also positively correlated themselves, with residential and commercial development, natural systems modifications, invasive species, and pollution having positive associations. Conservation research is aiming to further quantify the effects of threats, such as pollution, on carnivorous plants, as well as to quantify the extinction risks. Only 17% of species had been assessed as of 2011, according to the IUCN. Carnivorous plant conservation will help maintain important ecosystems and prevent secondary extinctions of specialist species that rely on them such as foundation species which may seek refuge or rely on certain plants for their existence. Research suggests a holistic approach, targeted at the habitat-level of carnivorous plants, may be required for successful conservation.
Classification
The classification of all flowering plants is currently in a state of flux. In the Cronquist system, the Droseraceae and Nepenthaceae were placed in the order Nepenthales, based on the radial symmetry of their flowers and their possession of insect traps. The Sarraceniaceae was placed either in the Nepenthales, or in its own order, the Sarraceniales. The Byblidaceae, Cephalotaceae, and Roridulaceae were placed in the Saxifragales; and the Lentibulariaceae in the Scrophulariales (now subsumed into the Lamiales).
In more modern classification, such as that of the Angiosperm Phylogeny Group, the families have been retained, but they have been redistributed amongst several disparate orders. It is also recommended that Drosophyllum be considered in a monotypic family outside the rest of the Droseraceae, probably more closely allied to the Dioncophyllaceae. The current recommendations are shown below (only carnivorous genera are listed):
Dicots
Asterales (sunflower and daisy order)
Stylidiaceae
Stylidium (trigger plants, a borderline carnivore)
Caryophyllales, (carnation order)
Dioncophyllaceae
Triphyophyllum (a tropical liana)
Drosophyllaceae
Drosophyllum (Portuguese dewy pine)
Droseraceae (sundew family)
Aldrovanda (waterwheel plant)
Dionaea (Venus flytrap)
Drosera (sundews)
†Droserapollis
†Droserapites
†Droseridites
†Fischeripollis
†Saxonipollis
Nepenthaceae (tropical pitcher-plant family)
Nepenthes (tropical pitcher plants or monkey-cups, including Anurosperma)
Ericales (heather order)
Roridulaceae
Roridula (a borderline carnivore)
Sarraceniaceae (trumpet pitcher family)
Sarracenia (North American trumpet pitchers)
Darlingtonia (cobra plant/lily)
Heliamphora (sun or marsh pitchers)
Lamiales (mint order)
Byblidaceae
Byblis (rainbow plants)
Lentibulariaceae (bladderwort family)
Pinguicula (butterworts)
Genlisea (corkscrew plant)
Utricularia (bladderworts, including Polypompholyx, the fairy aprons or pink petticoats and Biovularia an obsolete genus)
Martyniaceae (all borderline carnivores, related to the sesame plant)
Ibicella
Plantaginaceae (plantain family)
Philcoxia (recently discovered carnivorous genus feeding on nematodes).
Oxalidales (wood sorrel order)
Cephalotus (Albany pitcher plant)
Monocots
Alismatales (water plantain order)
Tofieldiaceae
Triantha occidentalis
Poales (grass order)
Bromeliaceae (bromeliad or pineapple family)
Brocchinia (a terrestrial bromeliad)
Catopsis (a borderline carnivore)
Eriocaulaceae (pipewort family)
Paepalanthus bromelioides (a borderline carnivore)
Gallery of prey
Cultivation
In horticulture, carnivorous plants are considered a curiosity or a rarity, but are becoming more common in cultivation with the advent of mass-production tissue-culture propagation techniques. Venus flytraps are still the most commonly grown, usually available at garden centers and hardware stores, sometimes offered alongside other easy-to-grow varieties. Nurseries that specialise in growing carnivorous plants exclusively also exist, more uncommon or demanding varieties of carnivorous plants can be obtained from specialist nurseries. California Carnivores is a notable example of such a nursery in the US that specialises in the cultivation of carnivorous plants. It is owned and operated by horticulturalist Peter D'Amato. Rob Cantley's Borneo Exotics in Sri Lanka is a large nursery that sells worldwide.
Although different species of carnivorous plants have different cultivation requirements in terms of sunlight, humidity, soil moisture, etc., there are commonalities. Most carnivorous plants require rainwater, or water that has been distilled or deionised by reverse osmosis. Common tap or drinking water contains minerals (particularly calcium salts) that will quickly build up and kill the plant. This is because most carnivorous plants have evolved in nutrient-poor, acidic soils and are consequently extreme calcifuges. They are therefore very sensitive to excessive soil-borne nutrients. Since most of these plants are found in bogs, almost all are very intolerant of drying. There are exceptions:
tuberous sundews require a dry (summer) dormancy period, and Drosophyllum requires much drier conditions than most.
Outdoor-grown carnivorous plants generally catch more than enough insects to keep themselves properly fed. Insects may be fed to the plants by hand to supplement their diet; however, carnivorous plants are generally unable to digest large non-insect food items; bits of hamburger, for example, will simply rot, and this may cause the trap, or even the whole plant, to die.
A carnivorous plant that catches no insects at all will rarely die, although its growth may be impaired. In general, these plants are best left to their own devices: after underwatering with tap-water, the most common cause of Venus flytrap death is prodding the traps to watch them close and feeding them inappropriate items.
Most carnivorous plants require bright light, and most will look better under such conditions, as this encourages them to synthesise red and purple anthocyanin pigments, (or betalain pigments within Caryophyllales). Nepenthes and Pinguicula will do better out of full sun, but most other species are happy in direct sunlight.
Carnivores mostly live in bogs, and those that do not are generally tropical. Hence, most require high humidity. On a small scale, this can be achieved by placing the plant in a wide saucer containing pebbles that are kept permanently wet. Small Nepenthes species grow well in large terraria.
Many carnivores are native to cold temperate regions and can be grown outside in a bog garden year-round. Most Sarracenia can tolerate temperatures well below freezing, despite most species being native to the southeastern United States. Species of Drosera and Pinguicula also tolerate subfreezing temperatures. Nepenthes species, which are tropical, require temperatures from 20 to 30 °C (70 to 90°F) to thrive.
Carnivorous plants require appropriate nutrient-poor soil. Most appreciate a 3:1 mixture of Sphagnum peat to sharp horticultural sand (coir is an acceptable, and more ecofriendly substitute for peat). Nepenthes will grow in orchid compost or in pure Sphagnum moss.
Carnivorous plants are themselves susceptible to infestation by parasites such as aphids or mealybugs. Although small infestations can be removed by hand, larger infestations necessitate use of an insecticide.
Isopropyl alcohol (rubbing alcohol) is effective as a topical insecticide, particularly on scale insects. Diazinon is an excellent systemic insecticide that is tolerated by most carnivorous plants. Malathion and Acephate (Orthene) have also been reported as tolerable by carnivorous plants.
Although insects can be a problem, by far the biggest killer of carnivorous plants (besides human maltreatment) is grey mold (Botrytis cinerea). This thrives under warm, humid conditions and can be a real problem in winter. To some extent, temperate carnivorous plants can be protected from this pathogen by ensuring that they are kept cool and well ventilated in winter and that any dead leaves are removed promptly. If this fails, a fungicide is in order.
The easiest carnivorous plants for beginners are those from the cool temperate zone. These plants will do well under cool greenhouse conditions (minimum 5 °C; 40°F in winter, maximum 25 °C; 75°F in summer) if kept in wide trays of acidified or rain water during summer and kept moist during winter:
Drosera capensis, the Cape sundew: attractive strap-leaved sundew, pink flowers, very tolerant of maltreatment.
Drosera binata, the fork-leaved sundew: large, Y-shaped leaves.
Sarracenia flava, the yellow trumpet pitcher: yellow, attractively veined leaves, yellow flowers in spring.
Pinguicula grandiflora, the common butterwort: purple flowers in spring, hibernates as a bud (hibernaculum) in winter. Fully hardy.
Pinguicula moranensis, the Mexican butterwort: pink flowers, non-carnivorous leaves in winter.
Venus flytraps will do well under these conditions but are actually rather difficult to grow: even if treated well, they will often succumb to grey mold in winter unless well ventilated. Some of the lowland Nepenthes are very easy to grow as long as they are provided with relatively constant, hot and humid conditions.
Medicinal uses
A study published in 2009 by researchers from Tel Aviv University indicates that secretions produced by carnivorous plants contain compounds that have anti-fungal properties and may lead to the development of a new class of anti-fungal drugs that will be effective against infections that are resistant to current anti-fungal drugs.
Cultural depictions
In 1789, Erasmus Darwin described Drosera in the second part of his poem The Botanic Garden:
However, Erasmus Darwin and others of his generation assumed that the "wonderful contrivance[s]" of carnivorous plants were solely defense mechanisms to "prevent various insects from plundering the honey, or devouring the seed". They realized that the plants were killing insects, but did not understand why. Erasmus Darwin's grandson, Charles Darwin, and great-grandson, Francis Darwin, spent many years studying carnivorous plants. Charles Darwin recognized and described the significance of plant carnivory for nutrition.
In 1860, residents of Providence, Rhode Island, dug up the grave of that state's founder Roger Williams, intending to move his remains to a new memorial in his honor. They found only teeth, nails, bone fragments, and an apple tree root that had grown along where his body had been, forking midway to follow his legs. The now-severed root forms its own sort of memorial, and has been called "The tree (or root) that ate Roger Williams".
Possibly the earliest published account of a man-eating plant was a literary fabrication that first appeared in 1874. The story of Crinoida dajeeana, also known as the Devil Tree of Madagascar or Man-Eating Tree of Madagascar, first appeared in the daily edition of the New York World on 26 April 1874, and again in the weekly edition two days later. It purported to be from a German explorer named "Karl Leche" (also spelled as Karl or Carl Liche in later accounts), who described seeing a woman fed to a tree as a sacrifice by the "little known but cruel" "Mkodo tribe" of Madagascar.
Authorship of the fantastical story would later be attributed by Frederick Maxwell Somers to one Edmund Spencer in the August 1888 issue of the magazine Current Literature. The story was reprinted widely, appearing as far away as the South Australian Register in 1881, where it was accompanied by an illustration of a tree consuming a woman. The account has been debunked as pure myth, and Dr. Liche, the Mkodos, and the tree itself were all fabrications.
Crinoida dajeeana notwithstanding, carnivorous plants are credited with widely entering the popular imagination through the nonfiction publications of Charles Darwin. Insectivorous Plants (1875), followed by The Power of Movement in Plants (1880), challenged the idea of a what a plant was and what it was capable of doing, and inspired authors like Arthur Conan Doyle to imagine enormous and sometimes mobile man-eaters. Doyle modeled the sticky end of a character in "The American's Tale" (1880) on a venus flytrap. H. G. Wells imagined a tentacular blood-sucking plant in "The Flowering of the Strange Orchid" (1894).
Since then, carnivorous plants have been the subject of popular interest and exposition, much of it highly inaccurate. Typically, these fictional depictions include exaggerated characteristics, such as enormous size or possession of abilities beyond the realm of reality, and can be viewed as a kind of artistic license.
In a 1939 pamphlet on carnivorous plants written for the Field Museum, Sophia Prior recounts the Man-Eating Tree of Madagascar and other "stories of vegetable monsters". She dismisses them all as fables, and notes that they are invariably set in locales that are "indefinite" and "difficult of access".
Fictional carnivorous plants have been featured in books, movies, television series, and video games. Some, such as the mockumentary The Hellstrom Chronicle (1971), use accurate depictions of carnivorous plants for cinematic purposes, while others depend more heavily on imagination.
Two of the most famous examples of fictional carnivorous plants in popular culture are the triffids of John Wyndham's 1951 novel The Day of the Triffids and Audrey Jr./II, the man-eating plant in the 1960s black comedy The Little Shop of Horrors and its subsequent stage musical adaptation.
| Biology and health sciences | Botany | null |
18952748 | https://en.wikipedia.org/wiki/Cuticle | Cuticle | A cuticle (), or cuticula, is any of a variety of tough but flexible, non-mineral outer coverings of an organism, or parts of an organism, that provide protection. Various types of "cuticle" are non-homologous, differing in their origin, structure, function, and chemical composition.
Human anatomy
In human anatomy, "cuticle" can refer to several structures, but it is used in general parlance, and even by medical professionals, to refer to the thickened layer of skin surrounding fingernails and toenails (the eponychium), and to refer to the superficial layer of overlapping cells covering the hair shaft (cuticula pili), consisting of dead cells, that locks the hair into its follicle. It can also be used as a synonym for the epidermis, the outer layer of skin.
Cuticle of invertebrates
In zoology, the invertebrate cuticle or cuticula is a multi-layered structure outside the epidermis of many invertebrates, notably arthropods and roundworms, in which it forms an exoskeleton (see arthropod exoskeleton).
The main structural components of the nematode cuticle are proteins, highly cross-linked collagens and specialised insoluble proteins known as "cuticlins", together with glycoproteins and lipids.
The main structural component of arthropod cuticle is chitin, a polysaccharide composed of N-acetylglucosamine units, together with proteins and lipids. The proteins and chitin are cross-linked. The rigidity is a function of the types of proteins and the quantity of chitin. It is believed that the epidermal cells produce protein and also monitors the timing and amount of protein to be incorporated into the cuticle.
Often, in the cuticle of arthropods, structural coloration is observed, produced by nanostructures. In the mealworm beetle, Tenebrio molitor, cuticular color may suggest pathogen resistance in that darker individuals are more resistant to pathogens compared to more tan individuals.
Botany
In botany, plant cuticles are protective, hydrophobic, waxy coverings produced by the epidermal cells of leaves, young shoots and all other aerial plant organs. Cuticles minimize water loss and effectively reduce pathogen entry due to their waxy secretion. The main structural components of plant cuticles are the unique polymers cutin or cutan, impregnated with wax. Plant cuticles function as permeability barriers for water and water-soluble materials. They prevent plant surfaces from becoming wet and also help to prevent plants from drying out. Xerophytic plants such as cacti have very thick cuticles to help them survive in their arid climates. Plants that live in range of sea's spray also may have thicker cuticles that protect them from the toxic effects of salt.
Some plants, particularly those adapted to life in damp or aquatic environments, have an extreme resistance to wetting. A well-known example is the sacred lotus. This adaptation is not purely the physical and chemical effect of a waxy coating but depends largely on the microscopic shape of the surface. When a hydrophobic surface is sculpted into microscopic, regular, elevated areas, sometimes in fractal patterns, too high and too closely spaced for the surface tension of the liquid to permit any flow into the space between the plateaus, then the area of contact between liquid and solid surfaces may be reduced to a small fraction of what a smooth surface might permit. The effect is to reduce wetting of the surface substantially.
Structural coloration is also observed in the cuticles of plants (see, as an example, the so-called "marble berry", Pollia condensata.
Mycology
"Cuticle" is one term used for the outer layer of tissue of a mushroom's basidiocarp, or "fruit body". The alternative term "pileipellis", Latin for "skin" of a "cap" (meaning "mushroom") might be technically preferable, but is perhaps too cumbersome for popular use. It is the part removed in "peeling" mushrooms. On the other hand, some morphological terminology in mycology makes finer distinctions, such as described in the article on the "pileipellis". Be that as it may, the pileipellis (or "peel") is distinct from the trama, the inner fleshy tissue of a mushroom or similar fruiting body, and also from the spore-bearing tissue layer, the hymenium.
| Biology and health sciences | Animal ontogeny | null |
18952853 | https://en.wikipedia.org/wiki/Agave | Agave | Agave (; ; ) is a genus of monocots native to the arid regions of the Americas. The genus is primarily known for its succulent and xerophytic species that typically form large rosettes of strong, fleshy leaves.
Many plants in this genus may be considered perennial, because they require several to many years to mature and flower. However, most Agave species are more accurately described as monocarpic rosettes or multiannuals, since each individual rosette flowers only once and then dies; a small number of Agave species are polycarpic.
Along with plants from the closely related genera Yucca, Hesperoyucca, and Hesperaloe, various Agave species are popular ornamental plants in hot, dry climates, as they require very little supplemental water to survive. Most Agave species grow very slowly. Some Agave species are known by the common name "century plant".
is a Spanish word that refers to all of the large-leafed plants in the Asparagaceae family, including agaves and yuccas. Maguey flowers are eaten in many indigenous culinary traditions of Mesoamerica.
Description
The succulent leaves of most Agave species have sharp marginal teeth, an extremely sharp terminal spine, and are very fibrous inside. The stout stem is usually extremely short, which may make the plant appear as though it is stemless.
Agave rosettes are mostly monocarpic, though some species are polycarpic. During flowering, a tall stem or "mast" ("quiote" in Mexico), which can grow to be high, grows apically from the center of the rosette and bears a large number of short, tubular flowers and sometimes vegetatively produced bulbils (a form of asexual reproduction). After pollination/fertilization and subsequent fruit development, in monocarpic species, the original rosette dies. However, throughout the lifetime of many Agave species, rhizomatous suckers develop above the roots at the base of the rosette. These suckers go on to form new plants after the original rosette desiccates and dies. Not all agaves produce suckers throughout their lifetimes; some species rarely or never produce suckers, while others may only develop suckers after final maturation with inflorescence. Some varieties can live for 60 years before flowering.
Agaves can be confused with cacti, aloes, or stonecrops, but although these plants all share similar morphological adaptations to arid environments (e.g. succulence), each group belongs to a different plant family and probably experienced convergent evolution. Further, cactus (Cactaceae) and stonecrop (Crassulaceae) lineages are eudicots, while aloes (Asphodelaceae) and agaves (Asparagaceae) are monocots.
Adaptations
The agave root system, consisting of a network of shallow rhizomes, allows the agave to efficiently capture moisture from rain, condensation, and dew. In addition to growing from seeds, most agaves produce 'pups' – young plants from runners. Agave vilmoriniana (the octopus agave) produces hundreds of pups on its bloom stalk. Agave leaves store the plant's water and are crucial to its continued existence. The coated leaf surface prevents evaporation. The leaves also have sharp, spiked edges. The spikes discourage predators from eating the plant or using it as a source of water and are so tough that ancient peoples used them for sewing needles. The sap is acidic. Some agaves bloom at a height up to so that they are far out of reach to animals that might attack them. Smaller species, such as Agave lechuguilla, have smaller bloom stalks.
Taxonomy
The genus name Agave come from the Ancient Greek from agauós meaning "illustrious, noble" having to do with very tall flower spikes found on its many species.
The genus Agave was erected by Carl Linnaeus in 1753, initially with four species. The first listed was Agave americana, now the type species. In the Cronquist system and others, Agave was placed in the family Liliaceae, but phylogenetic analyses of DNA sequences later showed it did not belong there. In the APG II system, Agave was placed in the segregated family Agavaceae. When this system was superseded by the APG III system in 2009, the Agavaceae were subsumed into the expanded family Asparagaceae, and Agave was treated as one of 18 genera in the subfamily Agavoideae, a position retained in the APG IV system of 2016.
Agaves and close relatives have long presented significant taxonomic difficulty. These difficulties could be due to the relatively young evolutionary age of the group (major diversification events of the group most likely occurred 8–10 million years ago), ease of hybridization between species (and even genera), incomplete lineage sorting, and long generation times. Within a species, morphological variations can be considerable, especially in cultivation; a number of named species may actually just be variants of original wild-type species that horticulturalists bred to appear unique in cultivation.
Commonly grown species
Some commonly grown species include Agave americana, A. angustifolia, A. attenuata, A. murpheyi, A. palmeri, A. parryi, A. parviflora, A. tequilana, A. victoriae-reginae, and A. vilmoriniana.
A. americana
One of the most familiar species is A. americana, a native of tropical America. Common names include century plant, maguey (in Mexico), or American aloe (though not related to the genus Aloe). The name "century plant" refers to the long time the plant takes to flower. The number of years before flowering occurs depends on the vigor of the individual plant, the richness of the soil, and the climate; during these years, the plant is storing in its fleshy leaves the nourishment required for the effort of flowering.
A. americana, century plant, was introduced into southern Europe about the middle of the 16th century and is now naturalized as well as widely cultivated as an ornamental, as it is in the Americas. In the variegated forms, the leaf has a white or yellow marginal or central stripe. As the leaves unfold from the center of the rosette, the impression of the marginal spines is conspicuous on the still erect younger leaves. The plant is reported being hardy to −9.5 to −6.5 °C or Zone 8b 15-20f. Being succulents, they tend to rot if kept too wet. In areas such as America's Pacific Northwest, they might be hardy for cold winter temperatures, but need protection from winter rain. They mature very slowly and die after flowering but are easily propagated by the offsets from the base of the stem.
A. americana (a blue variety) occurs in abundance in the Karoo, and arid highland regions of South Africa. Introduced by the British settlers in 1820, the plant was originally cultivated and used as emergency feed for livestock. Today, it is used mainly for the production of syrup and sugar.
A. attenuata
A. attenuata is a native of central Mexico and is uncommon in its natural habitat. Unlike most species of agave, A. attenuata has a curved flower spike from which it derives one of its numerous common names – the foxtail agave. It is also commonly grown as a garden plant. Unlike many agaves, A. attenuata has no teeth or terminal spines, making it an ideal plant for areas adjacent to footpaths. Like all agaves, it is a succulent and requires little water or maintenance once established.
A. tequilana
Agave azul (blue agave) is used in the production of tequila. It is native to the Caribbean as well as many regions of Mexico like Colima, Nayarit, Jalisco and more. In 2001, the Mexican government and European Union agreed upon the classification of tequila and its categories. All 100% blue agave tequila must be made from the A. tequilana 'Weber's Blue' agave plant, to rigorous specifications and only in the state of Jalisco. Blue agave is significantly different from other types of agave because it is higher in fructose and much sweeter compared to the rest. It is also the primary source for agave syrup, a nectary sweetener made for consumption.
Ecology
Agave species are used as food plants by the larvae of some Lepidoptera (butterfly and moth) species, including Batrachedra striolata, which has been recorded on A. shawii.
Toxicity
Some species contain components in their juice which can cause dermatitis for some people.
Uses
The ethnobotany of the agave was described by William H. Prescott in 1843:
But the miracle of nature was the great Mexican aloe, or maguey, whose clustering pyramids of flowers, towering above their dark coronals of leaves, were seen sprinkled over many a broad acre of the table-land. As we have already noticed its bruised leaves afforded a paste from which paper was manufactured, its juice was fermented into an intoxicating beverage, pulque, of which the natives, to this day, are extremely fond; its leaves further supplied an impenetrable thatch for the more humble dwellings; thread, of which coarse stuffs were made, and strong cords, were drawn from its tough and twisted fibers; pins and needles were made from the thorns at the extremity of its leaves; and the root, when properly cooked, was converted into a palatable and nutritious food. The agave, in short, was meat, drink, clothing, and writing materials for the Aztec! Surely, never did Nature enclose in so compact a form so many of the elements of human comfort and civilization!
The four major edible parts of the agave are the flowers, the leaves, the stalks or basal rosettes, and the sap (in Spanish: aguamiel, meaning "honey water"). The sap of some species can also be used as soap.
Food and fiber
Each agave plant produces several pounds of edible flowers during its final season. The stalks, which are ready during the summer, before the blossom, weigh several pounds each. Roasted, they are sweet and can be chewed to extract the sap or aguamiel, like sugarcane. When dried out, the stalks can be used to make didgeridoos. The leaves may be collected in winter and spring, when the plants are rich in sap, for eating. The leaves of several species also yield fiber, for instance, A. sisalana, the sisal hemp, and A. decipiens, the false sisal hemp. A. americana is the source of pita fiber, and is used as a fiber plant in Mexico, the West Indies, and southern Europe.
The agave, especially A. murpheyi, was a major food source for the prehistoric indigenous people of the Southwestern United States. The Hohokam of southern Arizona cultivated large areas of agave. In southern California and the Baja California Peninsula, the roasted hearts of A. shawii and A. deserti were historically among the most important foods for the Cahuilla, Kumeyaay, Kiliwa, and Paipai peoples, leaving ubiquitous archeological evidence in the form of agave-roasting pits throughout the region.
The Navajo similarly found many uses for the agave plant. A beverage is squeezed from the baked fibers, and the heads can be baked or boiled, pounded into flat sheets, sun dried, and stored for future use. The baked, dried heads are also boiled and made into an edible paste, eaten whole, or made into soup. The leaves are eaten boiled, and the young, tender flowering stalks and shoots are roasted and eaten as well. The fibers are used to make rope, the leaves are used to line baking pits, and the sharp-pointed leaf tips are used to make basketry awls.
During the development of the inflorescence, sap rushes to the base of the young flower stalk. Agave syrup (commonly called agave nectar), a sweetener derived from the sap, is used as an alternative to sugar in cooking, and can be added to breakfast cereals as a binding agent. Extracts from agave leaves are under preliminary research for their potential use as food additives.
Beverages and tequila
The sap of A. americana and other species is used in Mexico and Mesoamerica to produce pulque, an alcoholic beverage. The flower shoot is cut out and the sap collected and subsequently fermented. By distillation, a spirit called mezcal is prepared; one of the best-known forms of mezcal is tequila. A. tequilana or A. tequilana var. azul is used in the production of tequila. A. angustifolia is widely used in the production of mezcal and pulque, though at least 10 other Agave species are also known to be used for this.
Research
Agave can be used as the raw material for industrial production of fructans as a prebiotic dietary fiber. Agave contains fructooligosaccharides, which are naturally occurring oligosaccharides that support safely subjecting peanut-allergic people to allergen immunotherapy. Resulting from its natural habitat in stressful environments, agave is under preliminary research for its potential use in germplasm conservation and in biotechnology to better anticipate the economic effects of global climate change. It may also have use as a bioethanol or bioenergy feedstock.
Gallery of species and cultivars
| Biology and health sciences | Monocots | null |
18952860 | https://en.wikipedia.org/wiki/Transpiration | Transpiration | Transpiration is the process of water movement through a plant and its evaporation from aerial parts, such as leaves, stems and flowers. It is a passive process that requires no energy expense by the plant. Transpiration also cools plants, changes osmotic pressure of cells, and enables mass flow of mineral nutrients. When water uptake by the roots is less than the water lost to the atmosphere by evaporation, plants close small pores called stomata to decrease water loss, which slows down nutrient uptake and decreases CO2 absorption from the atmosphere limiting metabolic processes, photosynthesis, and growth.
Water and nutrient uptake
Water is necessary for plants, but only a small amount of water taken up by the roots is used for growth and metabolism. The remaining 97–99.5% is lost by transpiration and guttation. Water with any dissolved mineral nutrients is absorbed into the roots by osmosis, which travels through the xylem by way of water molecule adhesion and cohesion to the foliage and out small pores called stomata (singular "stoma"). The stomata are bordered by guard cells and their stomatal accessory cells (together known as stomatal complex) that open and close the pore. The cohesion-tension theory explains how leaves pull water through the xylem. Water molecules stick together or exhibit cohesion. As a water molecule evaporates from the leaf's surface it pulls on the adjacent water molecule, creating a continuous water flow through the plant.
Two major factors influence the rate of water flow from the soil to the roots: the hydraulic conductivity of the soil and the magnitude of the pressure gradient through the soil. Both of these factors influence the rate of bulk flow of water moving from the roots to the stomatal pores in the leaves via the xylem. Mass flow of liquid water from the roots to the leaves is driven in part by capillary action, but primarily driven by water potential differences. If the water potential in the ambient air is lower than that in the leaf airspace of the stomatal pore, water vapor will travel down the gradient and move from the leaf airspace to the atmosphere. This movement lowers the water potential in the leaf airspace and causes evaporation of liquid water from the mesophyll cell walls. This evaporation increases the tension on the water menisci in the cell walls and decreases their radius, thus exerting tension in the cells' water. Because of the cohesive properties of water, the tension travels through the leaf cells to the leaf and stem xylem, where a momentary negative pressure is created as water is pulled up the xylem from the roots. In taller plants and trees, the force of gravity pulling the water inside can only be overcome by the decrease in hydrostatic pressure in the upper parts of the plants due to the diffusion of water out of stomata into the atmosphere.
Etymology
The word transpiration comes from the words trans, a Latin preposition that means "across," and spiration, which comes from the Latin verb spīrāre, meaning "to breathe." The motion suffix adds the meaning "the act of," creating the meaning, "the ACT of breathing across."
Capillary action
Capillary action is the process of a liquid flowing in narrow spaces without the assistance of, or even in opposition to, external forces like gravity. The effect can be seen in the drawing up of liquids between the hairs of a paint-brush, in a thin tube, in porous materials such as paper and plaster, in some non-porous materials such as sand and liquefied carbon fiber, or in a biological cell. It occurs because of intermolecular forces between the liquid and surrounding solid surfaces. If the diameter of the tube is sufficiently small, then the combination of surface tension (which is caused by cohesion within the liquid) and adhesive forces between the liquid and container wall act to propel the liquid.
Regulation
Plants regulate the rate of transpiration by controlling the size of the stomatal apertures. The rate of transpiration is also influenced by the evaporative demand of the atmosphere surrounding the leaf such as boundary layer conductance, humidity, temperature, wind, and incident sunlight. Along with above-ground factors, soil temperature and moisture can influence stomatal opening, and thus transpiration rate. The amount of water lost by a plant also depends on its size and the amount of water absorbed at the roots. Factors that effect root absorption of water include: moisture content of the soil, excessive soil fertility or salt content, poorly developed root systems, and those impacted by pathogenic bacteria and fungi such as pythium or rhizoctonia.
During a growing season, a leaf will transpire many times more water than its own weight. An acre of corn gives off about 3,000–4,000 gallons (11,400–15,100 liters) of water each day, and a large oak tree can transpire 40,000 gallons (151,000 liters) per year. The transpiration ratio is the ratio of the mass of water transpired to the mass of dry matter produced; the transpiration ratio of crops tends to fall between 200 and 1000 (i.e., crop plants transpire 200 to 1000 kg of water for every kg of dry matter produced).
Transpiration rates of plants can be measured by a number of techniques, including potometers, lysimeters, porometers, photosynthesis systems and thermometric sap flow sensors. Isotope measurements indicate transpiration is the larger component of evapotranspiration. Recent evidence from a global study of water stable isotopes shows that transpired water is isotopically different from groundwater and streams. This suggests that soil water is not as well mixed as widely assumed.
Desert plants have specially adapted structures, such as thick cuticles, reduced leaf areas, sunken stomata and hairs to reduce transpiration and conserve water. Many cacti conduct photosynthesis in succulent stems, rather than leaves, so the surface area of the shoot is very low. Many desert plants have a special type of photosynthesis, termed crassulacean acid metabolism or CAM photosynthesis, in which the stomata are closed during the day and open at night when transpiration will be lower.
Cavitation
To maintain the pressure gradient necessary for a plant to remain healthy they must continuously uptake water with their roots. They need to be able to meet the demands of water lost due to transpiration. If a plant is incapable of bringing in enough water to remain in equilibrium with transpiration an event known as cavitation occurs. Cavitation is when the plant cannot supply its xylem with adequate water so instead of being filled with water the xylem begins to be filled with water vapor. These particles of water vapor come together and form blockages within the xylem of the plant. This prevents the plant from being able to transport water throughout its vascular system. There is no apparent pattern of where cavitation occurs throughout the plant's xylem. If not effectively taken care of, cavitation can cause a plant to reach its permanent wilting point, and die. Therefore, the plant must have a method by which to remove this cavitation blockage, or it must create a new connection of vascular tissue throughout the plant. The plant does this by closing its stomates overnight, which halts the flow of transpiration. This then allows for the roots to generate over 0.05 mPa of pressure, and that is capable of destroying the blockage and refilling the xylem with water, reconnecting the vascular system. If a plant is unable to generate enough pressure to eradicate the blockage it must prevent the blockage from spreading with the use of pit pears and then create new xylem that can re-connect the vascular system of the plant.
Scientists have begun using magnetic resonance imaging (MRI) to monitor the internal status of the xylem during transpiration, in a non invasive manner. This method of imaging allows for scientists to visualize the movement of water throughout the entirety of the plant. It also is capable of viewing what phase the water is in while in the xylem, which makes it possible to visualize cavitation events. Scientists were able to see that over the course of 20 hours of sunlight more than 10 xylem vessels began filling with gas particles becoming cavitated. MRI technology also made it possible to view the process by which these xylem structures are repaired in the plant. After three hours in darkness it was seen that the vascular tissue was resupplied with liquid water. This was possible because in darkness the stomates of the plant are closed and transpiration no longer occurs. When transpiration is halted the cavitation bubbles are destroyed by the pressure generated by the roots. These observations suggest that MRIs are capable of monitoring the functional status of xylem and allows scientists to view cavitation events for the first time.
Effects on the environment
Cooling
Transpiration cools plants, as the evaporating water carries away heat energy due to its large latent heat of vaporization of 2260 kJ per liter.
| Biology and health sciences | Basics_3 | null |
18952932 | https://en.wikipedia.org/wiki/Hallucinogen | Hallucinogen | Hallucinogens are a large and diverse class of psychoactive drugs that can produce altered states of consciousness characterized by major alterations in thought, mood, and perception as well as other changes. Most hallucinogens can be categorized as either being psychedelics, dissociatives, or deliriants.
Etymology
The word hallucinogen is derived from the word hallucination. The term hallucinate dates back to around 1595–1605, and is derived from the Latin hallūcinātus, the past participle of (h)allūcināri, meaning "to wander in the mind."
Characteristics
Leo Hollister gave five criteria for classifying a drug as hallucinogenic. This definition is broad enough to include a wide range of drugs and has since been shown to encompass a number of categories of drugs with different pharmacological mechanisms and behavioral effects. Richard Glennon has thus given an additional two criteria that narrow the category down to classical hallucinogens. Hollister's criteria for hallucinogens were as follows:
in proportion to other effects, changes in thought, perception, and mood should predominate;
intellectual or memory impairment should be minimal;
stupor, narcosis, or excessive stimulation should not be an integral effect;
autonomic nervous system side effects should be minimal; and
addictive craving should be absent.
Glennon's additional criteria for classical hallucinogens are that the drugs in question must also:
bind at 5-HT2 serotonin receptors; and
be recognized by animals trained to discriminate the drug DOM from vehicle.
Nomenclature and taxonomy
Most hallucinogens can be categorized based on their pharmacological mechanisms as psychedelics (which are serotonergic), dissociatives (which are generally antiglutamatergic), or deliriants (which are generally anticholinergic). However, the pharmacological mechanisms of some hallucinogens, such as salvinorin A and ibogaine, do not fit into any of those categories. Entactogens and cannabinoids are also sometimes considered hallucinogens. Nonetheless, while the term hallucinogen is often used to refer to the broad class of drugs covered in this article, sometimes it is used to mean only classical hallucinogens (that is, psychedelics). Because of this, it is important to consult the definition given in a particular source. Because of the multi-faceted phenomenology brought on by hallucinogens, efforts to create standardized terminology for classifying them based on their subjective effects have not succeeded to date.
Classical hallucinogens or psychedelics have been described by many names. David E. Nichols wrote in 2004:
Robin Carhart-Harris and Guy Goodwin write that the term psychedelic is preferable to hallucinogen for describing classical psychedelics because of the term hallucinogens "arguably misleading emphasis on these compounds' hallucinogenic properties."
Certain hallucinogens are designer drugs, such as those in the 2C and 25-NB (NBOMe) families. A designer drug is a structural or functional analog of a controlled substance (hallucinogenic or otherwise) that has been designed to mimic the pharmacological effects of the original drug while at the same time avoid being classified as illegal (by specification as a research chemical) and/or avoid detection in standard drug tests.
Effects by type
Psychedelics (classical hallucinogens)
Despite several attempts that have been made, starting in the 19th and 20th centuries, to define common phenomenological structures (i.e., patterns of experience) brought on by classical psychedelics, a universally accepted taxonomy does not yet exist.
A prominent element of psychedelic experiences is visual alteration. Psychedelic visual alteration often includes spontaneous formation of complex flowing geometric visual patterning in the visual field. When the eyes are open, the visual alteration is overlaid onto the objects and spaces in the physical environment; when the eyes are closed the visual alteration is seen in the "inner world" behind the eyelids. These visual effects increase in complexity with higher dosages, and also when the eyes are closed. The visual alteration does not normally constitute hallucinations, because the person undergoing the experience can still distinguish between real and internally generated visual phenomena, though in some cases, true hallucinations are present. More rarely, psychedelic experiences can include complex hallucinations of objects, animals, people, or even whole landscapes.
A number of studies by Roland R. Griffiths and other researchers have concluded that high doses of psilocybin and other classic psychedelics trigger mystical experiences in most research participants. Mystical experiences have been measured by a number of psychometric scales, including the Hood Mysticism Scale, the Spiritual Transcendence Scale, and the Mystical Experience Questionnaire. The revised version of the Mystical Experience Questionnaire, for example, asks participants about four dimensions of their experience, namely the "mystical" quality, positive mood such as the experience of amazement, the loss of the usual sense of time and space, and the sense that the experience cannot be adequately conveyed through words. The questions on the "mystical" quality in turn probe multiple aspects: the sense of "pure" being, the sense of unity with one's surroundings, the sense that what one experienced was real, and the sense of sacredness. Some researchers have questioned the interpretation of the results from these studies and whether the framework and terminology of mysticism are appropriate in a scientific context, while other researchers have responded to those criticisms and argued descriptions of mystical experiences are compatible with a scientific worldview.
Link R. Swanson divides overarching scientific frameworks for understanding psychedelic experiences into two waves. In the first wave, encompassing nineteenth- and twentieth-century frameworks, he includes model psychosis theory (the psychotomimetic paradigm), filtration theory, and psychoanalytic theory. In the second wave of theories, encompassing twenty-first-century frameworks, Swanson includes entropic brain theory, integrated information theory, and predictive processing. It is from the paradigm of filtration theory that the term psychedelic derives. Aldous Huxley and Humphrey Osmond applied the pre-existing ideas of filtration theory, which held that the brain filters what enters into consciousness, to explain psychedelic experiences; Huxley believed that the brain was filtering reality itself and that psychedelics granted conscious access to "Mind at Large", whereas Osmond believed that the brain was filtering aspects of the mind out of consciousness. Swanson writes that Osmond's view seems "less radical, more compatible with materialist science, and less epistemically and ontologically committed" than Huxley's.
Dissociatives
Dissociatives produce analgesia, amnesia and catalepsy at anesthetic doses. They also produce a sense of detachment from the surrounding environment, hence "the state has been designated as dissociative anesthesia since the patient truly seems disassociated from his environment." Dissociative symptoms include the disruption or compartmentalization of "...the usually integrated functions of consciousness, memory, identity or perception."p. 523 Dissociation of sensory input can cause derealization, the perception of the outside world as being dream-like, vague or unreal. Other dissociative experiences include depersonalization, which includes feeling dissociated from one's personality; feeling unreal; feeling able to observe one's actions but not actively take control; being unable to associate with one's self in the mirror while maintaining rational awareness that the image in the mirror is the same person. In a 2004 paper, Daphne Simeon offered "...common descriptions of depersonalisation experiences: watching oneself from a distance (similar to watching a movie); candid out-of-body experiences; a sense of just going through the motions; one part of the self acting/participating while the other part is observing;...."
The classical dissociatives achieve their effect through blocking binding of the neurotransmitter glutamate to NMDA receptors (NMDA receptor antagonism) and include ketamine, methoxetamine (MXE), phencyclidine (PCP), dextromethorphan (DXM), and nitrous oxide. However, dissociation is also remarkably administered by salvinorin A's (the active constituent in Salvia divinorum shown to the left) potent κ-opioid receptor agonism, though usually described as a very atypical dissociative.
Some dissociatives can have CNS depressant effects, thereby carrying similar risks as opioids, which can slow breathing or heart rate to levels resulting in death (when using very high doses). DXM in higher doses can increase heart rate and blood pressure and still depress respiration. Inversely, PCP can have more unpredictable effects and has often been classified as a stimulant and a depressant in some texts along with being as a dissociative. While many have reported that they "feel no pain" while under the effects of PCP, DXM and Ketamine, this does not fall under the usual classification of anesthetics in recreational doses (anesthetic doses of DXM may be dangerous). Rather, true to their name, they process pain as a kind of "far away" sensation; pain, although present, becomes a disembodied experience and there is much less emotion associated with it. As for probably the most common dissociative, nitrous oxide, the principal risk seems to be due to oxygen deprivation. Injury from falling is also a danger, as nitrous oxide may cause sudden loss of consciousness, an effect of oxygen deprivation. Because of the high level of physical activity and relative imperviousness to pain induced by PCP, some deaths have been reported due to the release of myoglobin from ruptured muscle cells. High amounts of myoglobin can induce renal shutdown.
Many users of dissociatives have been concerned about the possibility of NMDA antagonist neurotoxicity (NAN). This concern is partly due to William E. White, the author of the DXM FAQ, who claimed that dissociatives definitely cause brain damage. The argument was criticized on the basis of lack of evidence and White retracted his claim. White's claims and the ensuing criticism surrounded original research by John Olney.
In 1989, John Olney discovered that neuronal vacuolation and other cytotoxic changes ("lesions") occurred in brains of rats administered NMDA antagonists, including PCP and ketamine. Repeated doses of NMDA antagonists led to cellular tolerance and hence continuous exposure to NMDA antagonists did not lead to cumulative neurotoxic effects. Antihistamines such as diphenhydramine, barbiturates and even diazepam have been found to prevent NAN. LSD and DOB have also been found to prevent NAN.
Deliriants
Deliriants, as their name implies, induce a state of delirium in the user, characterized by extreme confusion and an inability to control one's actions. They are called deliriants because their subjective effects are similar to the experiences of people with delirious fevers. The term was introduced by David F. Duncan and Robert S. Gold to distinguish these drugs from psychedelics and dissociatives, such as LSD and ketamine respectively, due to their primary effect of causing delirium, as opposed to the more lucid states produced by the other hallucinogens.
Despite the fully legal status of several common deliriant plants, deliriants are largely unpopular as recreational drugs due to the severe, generally unpleasant and often dangerous nature of the hallucinogenic effects produced.
Typical or classical deliriants are those which are anticholinergic, meaning they block the muscarinic acetylcholine receptors. Many of these compounds are produced naturally by plant genera belonging to the nightshade family Solanaceae, such as Datura, Brugmansia and Latua in the New World and Atropa, Hyoscyamus and Mandragora in the Old World. These tropane alkaloids are poisonous and can cause death due to tachycardia-induced heart failure and hyperthermia even in small doses. Additionally, over-the-counter antihistamines such as diphenhydramine (brand name Benadryl) and dimenhydrinate (brand name Dramamine) also have an anticholinergic effect.
Uncured tobacco is also a deliriant due to its intoxicatingly high levels of nicotine.
History of use
Traditional religious and shamanic use
Historically, hallucinogens have been commonly used in religious or shamanic rituals. In this context they are referred to as entheogens, and are used to facilitate healing, divination, communication with spirits, and coming-of-age ceremonies. Evidence exists for the use of entheogens in prehistoric times, as well as in numerous ancient cultures, including Ancient Egyptian, Mycenaean, Ancient Greek, Vedic, Maya, Inca and Aztec cultures. The Upper Amazon is home to the strongest extant entheogenic tradition; the Urarina of the Peruvian Amazon, for instance, continue to practice an elaborate system of ayahuasca shamanism, coupled with an animistic belief system.
Shamans consume hallucinogenic substances in order to induce a trance. Once in this trance, shamans believe that they are able to communicate with the spirit world, and can see what is causing their patients' illness. The Aguaruna of Peru believe that many illnesses are caused by the darts of sorcerers. Under the influence of yaji, a hallucinogenic drink, Aguaruna shamans try to discover and remove the darts from their patients.
In the 1970s, Frida G. Surawicz and Richard Banta published a review of two case studies where hallucinogenic drug use appeared to play a role in "delusions of being changed into a wolf" (sometimes referred to as "lycanthropy," or being a "werewolf"). They described a patient whose delusion was thought to be caused by an altered state of consciousness "brought on by LSD and strychnine and continued casual marijuana use." The review was published in the Canadian Psychiatric Association Journal. While both central cases described white male patients from contemporary Appalachia, Surawicz and Banta generalized their conclusions about a link between hallucinogens and "lycanthropy," based on historical accounts that reference myriad types of pharmacologically-similar drug-use alongside descriptions of "lycanthropes."
Early scientific investigations
In an 1860 book, the mycologist Mordecai Cubitt Cooke differentiated a class of drugs roughly corresponding to hallucinogens from opiates, and in 1924 the toxicologist Louis Lewin described hallucinogens in depth under the name phantastica. From the 1920s on, work in psychopharmacology and ethnobotany resulted in more detailed knowledge of various hallucinogens. In 1943, Albert Hofmann discovered the hallucinogenic properties of lysergic acid diethylamide (LSD), which raised the prospect of hallucinogens becoming more broadly available.
Hallucinogens after World War II
After World War II there was an explosion of interest in hallucinogenic drugs in psychiatry, owing mainly to the invention of LSD. Interest in the drugs tended to focus on either the potential for psychotherapeutic applications of the drugs (see psychedelic psychotherapy), or on the use of hallucinogens to produce a "controlled psychosis", in order to understand psychotic disorders such as schizophrenia. By 1951, more than 100 articles on LSD had appeared in medical journals, and by 1961, the number had increased to more than 1000 articles.
At the beginning of the 1950s, the existence of hallucinogenic drugs was virtually unknown to the general public in the West. However this soon changed as several influential figures were introduced to the hallucinogenic experience. Aldous Huxley's 1953 essay The Doors of Perception, describing his experiences with mescaline, and R. Gordon Wasson's 1957 Life magazine article ("Seeking the Magic Mushroom") brought the topic into the public limelight. In the early 1960s, counterculture icons such as Jerry Garcia, Timothy Leary, Allen Ginsberg and Ken Kesey advocated the drugs for their psychedelic effects, and a large subculture of psychedelic drug users was spawned. Psychedelic drugs played a major role in catalyzing the major social changes initiated in the 1960s. As a result of the growing popularity of LSD and disdain for the hippies with whom it was heavily associated, LSD was banned in the United States in 1967. This greatly reduced the clinical research about LSD, although limited experiments continued to take place, such as those conducted by Reese Jones in San Francisco.
As early as the 1960s, research into the medicinal properties of LSD was being conducted. "Savage et al. (1962) provided the earliest report of efficacy for a hallucinogen in OCD, where after two doses of LSD, a patient who suffered from depression and violent obsessive sexual thoughts experienced dramatic and permanent improvement (Nichols 2004: 164)."
Starting in the mid-20th century, psychedelic drugs have received extensive attention in the Western world. They have been and are being explored as potential therapeutic agents in treating alcoholism, and other forms of drug addiction.
Legal status and attitudes
In the United States, classical hallucinogens (psychedelics) are in the most strictly prohibited class of drugs, known as Schedule 1 drugs. This classification was created for drugs that meet the three following characteristics: 1) they have no currently accepted medical use, 2) there is a lack of safety for their use under medical supervision, and 3) they have a high potential for abuse. However, pharmacologist David E. Nichols argues that hallucinogens were placed in this class for political rather than scientific reasons. In 2006, Albert Hofmann, the chemist who discovered LSD, said he believed LSD could be valuable when used in a medical rather than recreational context, and said it should be regulated in the same way as morphine rather than more strictly.
The Netherlands previously allowed psilocybin mushrooms to be sold, but in October 2007 the Dutch government moved to ban their sale following several widely publicized incidents involving tourists. In November 2020, Oregon became the first U.S. state to both decriminalize psilocybin and legalize it for therapeutic use, after Ballot Measure 109 passed.
Effects
Relationship between long-term use and mental illness
No clear connection has been made between psychedelic drugs and organic brain damage. However, hallucinogen persisting perception disorder (HPPD) is a diagnosed condition wherein certain visual effects of drugs persist for a long time, sometimes permanently, although the underlying cause and pathology remains unclear.
A large epidemiological study in the U.S. found that other than personality disorders and other substance use disorders, lifetime hallucinogen use was not associated with other mental disorders, and that risk of developing a hallucinogen use disorder was very low.
A 2019 systematic review and meta-analysis by Murrie et al. found that the transition rate from a diagnosis of hallucinogen-induced psychosis to that of schizophrenia was 26% (CI 14%-43%), which was lower than cannabis-induced psychosis (34%) but higher than amphetamine (22%), opioid (12%), alcohol (10%) and sedative (9%) induced psychoses. Transition rates were not affected by sex, country of the study, hospital or community location, urban or rural setting, diagnostic methods, or duration of follow-up. In comparison, the transition rate for brief, atypical and not otherwise specified psychosis was found to be 36%.
Effects on the brain
Different classes of hallucinogens have different pharmacological mechanisms of action. Psychedelics are 5-HT2A receptor agonists (serotonin 2A receptor agonists).
LSD, mescaline, psilocybin, and PCP are drugs that cause hallucinations, which can alter a person's perception of reality. LSD, mescaline, and psilocybin cause their effects by initially disrupting the interaction of nerve cells and the neurotransmitter serotonin. It is distributed throughout the brain and spinal cord, where the serotonin system is involved with controlling of the behavioral, perceptual, and regulatory systems. This also includes mood, hunger, body temperature, sexual behavior, muscle control, and sensory perception. Certain hallucinogens, such as PCP, act through a glutamate receptor in the brain which is important for perception of pain, responses to the environment, and learning and memory.
Thus far, there have been no properly controlled research studies on the specific effects of these drugs on the human brain, but smaller studies have shown some of the documented effects associated with the use of hallucinogens.
Psychotomimetic paradigm
While early researchers believed certain hallucinogens mimicked the effects of schizophrenia, it has since been discovered that some hallucinogens resemble endogenous psychoses better than others. PCP and ketamine are known to better resemble endogenous psychoses because they reproduce both positive and negative symptoms of psychoses, while psilocybin and related hallucinogens typically produce effects resembling only the positive symptoms of schizophrenia. While the serotonergic psychedelics (LSD, psilocybin, mescaline, etc.) do produce subjective effects distinct from NMDA antagonist dissociatives (PCP, ketamine, dextrorphan), there is obvious overlap in the mental processes that these drugs affect and research has discovered that there is overlap in the mechanisms by which both types of psychedelics mimic psychotic symptoms. One double-blind study examining the differences between DMT and ketamine hypothesized that classically psychedelic drugs most resemble paranoid schizophrenia while dissociative drugs best mimicked catatonic subtypes or otherwise undifferentiated schizophrenia. The researchers stated that their findings supported the view that "a heterogeneous disorder like schizophrenia is unlikely to be modeled accurately by a single pharmacological agent."
Chemistry
Classical hallucinogens (psychedelics) can be divided into three main chemical classes: tryptamines (such as psilocin and DMT), ergolines (such as LSD), and phenethylamines (such as mescaline). Tryptamines closely resemble serotonin chemically.
| Biology and health sciences | General concepts_2 | Health |
18952991 | https://en.wikipedia.org/wiki/Gila%20monster | Gila monster | The Gila monster (Heloderma suspectum, ) is a species of venomous lizard native to the Southwestern United States and the northwestern Mexican state of Sonora. It is a heavy, slow-moving reptile, up to long, and it is the only venomous lizard native to the United States. Its venomous close relatives, the four beaded lizards (all former subspecies of Heloderma horridum) inhabit Mexico and Guatemala. The Gila monster is sluggish in nature, so it is not generally dangerous and very rarely poses a real threat to humans. However, it has a fearsome reputation and is sometimes killed despite the species being protected by state law in Arizona.
History
The name "Gila" refers to the Gila River Basin in the U.S. states of Arizona and New Mexico, where the Gila monster was once plentiful. Heloderma means "studded skin", from the Ancient Greek words (), "the head of a nail or stud", and (), "skin". The species epithet suspectum comes from the describer, paleontologist Edward Drinker Cope. At first, this new specimen of Heloderma was misidentified and considered to be a northern variation of the beaded lizard already known to live in Mexico. He suspected that the lizard might be venomous due to the grooves in the teeth.
The Gila monster is the largest extant lizard species native to North America north of the Mexican border. Its snout-to-vent length ranges from . The tail is about 20% of the body size, and the largest specimens may reach in total length. Body mass is typically in the range of . They appear strong in their body structure with a stout snout, massive head, and "little"-appearing eyes, which can be protected by a nictitating membrane.
The Gila monster has four close living relatives, all of which are beaded lizards. There are three species in Mexico: Heloderma exasperatum, Heloderma horridum and Heloderma alvarezi, as well as another species in Guatemala: Heloderma charlesbogerti.
The evolutionary history of the Helodermatidae may be traced back to the Cretaceous period (145 to 166 million years ago), when Gobiderma pulchrum and Estesia mongolensis were present. The genus Heloderma has existed since the Miocene, when H. texana lived. Fragments of osteoderms from the Gila monster have been found in Late Pleistocene (10,000 to 8,000 years ago) deposits near Las Vegas, Nevada. Because the helodermatids have remained relatively unchanged morphologically, they are occasionally regarded as living fossils. Although the Gila monster appears closely related to the monitor lizards (varanids) of Africa, Asia, and Australia, their wide geographical separation and distinct features indicate that Heloderma is better placed in a separate family.
Skin
The scales of the head, back, and tail contain little pearl-shaped bones (osteoderms) similar to those found in the beaded lizards from further south. The scales of the belly are free from osteoderms. Female Gila monsters go through a total shed lasting about 2 weeks before depositing their eggs. The dorsal part is often shed in one large piece. Adult males normally shed in smaller segments in August. The young seem to be in constant shed. Adults have more or less yellow to pink colors on a black surface. Hatchlings have a uniform, simple, and less colorful pattern. This drastically changes within the first 6 months of their lives. Hatchlings from the northern area of the species' distribution have a tendency to retain most of their juvenile pattern. It has also been seen that Gila monsters in areas with darker rocks and substrate will have darker colorations.
The heads of males are very often larger and more triangular-shaped than in females. The length of the tail of the two sexes is statistically very similar, so it does not help in differentiation of the sexes. Individuals with stout tail ends occur in both nature and under human breeding.
Distribution and habitat
The Gila monster is found in the Southwestern United States and Mexico, across a range including Sonora, Arizona, and parts of California, Nevada, Utah, and New Mexico. No records have been given from Baja California. They inhabit scrubland, succulent desert, and oak woodland, seeking shelter in burrows, thickets, and under rocks in locations with a favorable microclimate and adequate humidity. Gila monsters rely heavily on the use of shelters and spend much of their time dwelling there. Often times these shelters are in rocky areas in Navajo Sandstone and basaltic lava flows. Gila monsters depend on water resources and can be observed in puddles of water after a summer rain. They avoid living in open areas, such as flats and open grasslands.
Ecology
Gila monsters spend 90% of their lifetime underground in burrows or rocky shelters. They are active in the morning during the dry season (spring and early summer). The lizards move to different shelters every 4–5 days up to the beginning of the summer season. By doing so, they optimize for a suitable microhabitat. Later in the summer, they may be active on warm nights or after a thunderstorm. They maintain a surface body temperature of about . Close to , they are able to decrease their body temperature by up to 2 °C (3.6 °F) by an activated, limited evaporation via the cloaca. One study investigating a population of Gila monsters in southwestern Utah noted that the lizard's activity peaked from late April to mid June. The average distance traveled during their bouts of activity was , but on occasion some lizards would travel distances greater than . During the Gila monster's active season of approximately 90 days, only ten days were spent active. Gila monsters are slow sprinters, but they have relatively high endurance and maximal aerobic capacity (VO2 max) compared to other lizards. They are preyed upon by coyotes, badgers and raptors. Hatchlings are preyed on by snakes, such as kingsnakes (Lampropeltis sp.).
Among adaptations to a dry environment is a slow metabolism, allowing them to use less than half the amount of energy expected for lizards of their size. Gila monsters, and possibly also the Mexican beaded lizard, store water in their urinary bladder and reabsorb it across the bladder epithelium. Their tail is used for energy storage in the form of fat.
Diet
The Gila monster's diet consists of a variety of food items – small mammals (such as young rabbits, hares, mice, ground squirrels, and other rodents), small birds, snakes, lizards, frogs, insects, other invertebrates, carrion, and the eggs of birds, lizards, snakes, and tortoises. Three to four extensive meals in spring are claimed to give them enough energy for a whole season. They can store fat in their tails and therefore do not need to eat often. Nevertheless, they feed whenever they come across suitable prey. Young Gila monsters can swallow up to 50% of their body weight in a single meal. Adults may eat up to one third of their body weight in one meal.
The Gila monster uses its extremely acute sense of smell to locate prey. The strong, two-ended tipped tongue, which is pigmented in black-blue colors, picks up scent molecules to be transferred to the opening of the Jacobson organ around the middle of the upper mouth cavern.
Prey may be crushed to death if large, or eaten alive, most of the time head first, and helped down by muscular contractions and neck flexing. After food has been swallowed, the Gila monster may immediately resume tongue flicking and search behavior for identifying more prey such as eggs or young in nests. Gila monsters are able to climb trees, cacti, and even fairly straight, rough-surfaced walls.
Venom
Pioneer beliefs
In the Old West, the pioneers believed a number of myths about the Gila monster, including that the lizard had foul or toxic breath and that its bite was fatal. The Tombstone Epitaph of Tombstone, Arizona, wrote about a Gila monster that a local person caught on May 14, 1881:
On May 8, 1890, southeast of Tucson, Arizona Territory, Empire Ranch owner Walter Vail captured and thought he had killed a Gila monster. He tied it to his saddle and it bit the middle finger of his right hand and would not let go. A ranch hand pried open the lizard's mouth with a pocket knife, cut open his finger to stimulate bleeding, and then tied saddle strings around his finger and wrist. They summoned Dr. John C. Handy of Tucson, who took Vail back to Tucson for treatment, but Vail experienced swollen and bleeding glands in his throat for sometime afterward.
Dr. Handy's friend, Dr. George Goodfellow of Tombstone, was among the first to research the actual effects of Gila monster venom. Scientific American reported in 1890, "The breath is very fetid, and its odor can be detected at some little distance from the lizard. It is supposed that this is one way in which the monster catches the insects and small animals which form a part of its food supplythe foul gas overcoming them." Goodfellow offered to pay local residents $5.00 for Gila monster specimens. He bought several and collected more on his own. In 1891, he purposely provoked one of his captive lizards into biting him on his finger. The bite made him ill and he spent the next five days in bed, but he completely recovered. When Scientific American ran another ill-founded report on the lizard's ability to kill people, he wrote in reply and described his own studies and personal experience. He wrote that he knew several people who had been bitten by Gila monsters but had not died from the bite.
Venom delivery
The Gila monster produces venom in modified salivary glands at the end of its lower jaws, unlike snakes, whose venom is produced in glands behind the eyes. The Gila monster lacks strong musculature in glands above the eyes; instead, in Heloderma, the venom is propelled from the gland via a tubing to the base of the lower teeth and then by capillary forces into two grooves of the tooth and then chewed into the victim. The teeth are tightly anchored to the jaw (pleurodont). Broken and regular replacement teeth have to wait every time to go into position in a determinate "wavelike" sequence. They change/replace their teeth during their entire life. The Gila monster's bright colors might be suitable to teach predators not to bother this "painful" creature. Because the Gila monster's prey consists mainly of eggs, small animals, and otherwise "helpless" prey, the Gila monster's venom is thought to have evolved for defensive rather than for hunting use.
Toxicity
The venom of a Gila monster is normally not fatal to healthy adult humans. One fatality has been confirmed since 1930, on February 16, 2024, and the rare fatalities recorded before 1930 occurred in adults who were intoxicated by alcohol or had mismanaged the treatment of the bite. The Gila monster can bite quickly and may not release the victim without intervention. If bitten, the victim may attempt to fully submerge the lizard in water, pry the jaws open with a knife or stick, or physically yank the lizard free. While pulling the lizard directly increases risk of severe lacerations from the lizard's sharp teeth, it may also mitigate envenomation. Symptoms of the bite include excruciating pain, edema, and weakness associated with a rapid drop in blood pressure.
YouTuber Coyote Peterson described the bite as "like hot lava coursing through your veins" and claimed it was "the worst pain [he] had ever experienced". It is generally regarded as the most painful venom produced by any vertebrate.
More than a dozen peptides and other substances have been isolated from the Gila monster's venom, including hyaluronidase, serotonin, phospholipase A2, and several kallikrein-like glycoproteins responsible for the pain and edema caused by a bite, without producing a compartment syndrome. Four potentially lethal toxins have been isolated from the Gila monster's venom, which cause hemorrhage in internal organs and exophthalmos (bulging of the eyes), and helothermine, which causes lethargy, partial paralysis of the limbs, and hypothermia in rats. Some are similar in action of the vasoactive intestinal peptide (VIP), which relaxes smooth muscle and regulates water and electrolyte secretion between the small and large intestines. These bioactive peptides are able to bind to VIP receptors in many different human tissues. One of these, helodermin, has been shown to inhibit the growth of lung cancer.
In February 2024, a Colorado man died after sustaining a 4 minute long bite from a pet Gila Monster. Autopsy reports determined cause of death to come from the animal bite/venom.
Relationship to GLP-1 agonists
The constituents of H. suspectum venom that have received the most attention from researchers are the bioactive peptides, including helodermin, helospectin, exendin-3, and exendin-4. Exendin-4, which is specific for H. suspectum, has formed the basis of a class of medications for the treatment of type 2 diabetes and obesity/metabolic syndrome, known as the GLP-1 receptor agonists.
In 2005, the U.S. Food and Drug Administration approved the drug exenatide (marketed as Byetta) for the management of type 2 diabetes. It is a synthetic modification of the protein exendin-4, that was originally isolated from the Gila monster's venom. In a 3-year study with people with type 2 diabetes, exenatide showed healthy sustained glucose levels. The effectiveness is because the lizard protein is 53% identical to glucagon-like peptide-1 analog (GLP-1), a hormone released from the human digestive tract that helps to regulate insulin and glucagon. Using a sophisticated injection formula with sustained release of the drug, the lizard protein remains effective much longer than the human hormone. This helps diabetics keep their blood glucose levels under control for a week by a single injection. Exenatide also slows the emptying of the stomach and causes a decrease in appetite, contributing to weight loss.
Life cycle
The Gila monster emerges from brumation in early March. Gila monsters sexually mature at 4–5 years old. It mates in April and May. The male initiates courtship by flicking his tongue to search for the female's scent. If the female rejects his advances, she will bite him and chase him away. When successful, copulation has been observed in captivity to last from 15 minutes to two and a half hours. There is only a single record of attempted mating outside of a shelter. The female lays eggs at the end of May into June. A clutch may consist of up to eight (rarely more than six) eggs. The incubation in captivity lasts about 5 months, depending on the incubation temperature. The hatchlings are about long and can bite and inject venom as soon as they are hatched.
The egg development and hatching time of young in the wild has been a subject of ongoing speculation. The first model stated that youngsters hatch in fall and stay underground. The second theory postulated a nearly developed embryo remains inside the egg over winter and hatches in spring. Hatchlings (weight about ) are observed at the end of April to early June.
Discussions of the exact egg development and hatching cycle of the Gila monster came to an abrupt and unexpected end on October 28, 2016, when a backhoe was digging at the outer walls of a house in a suburb of northern Tucson. The backhoe extracted a nest of a female Gila monster with five eggs in the process of hatching. The Gila monster is now known to hatch near the end of October and immediately proceed into hibernation without surfacing. They then appear on the surface from May through June the following year when prey should be abundant.
In summer, Gila monsters gradually spend less time on the surface to avoid the hottest part of the season; occasionally, they may be active at night. Females that have laid eggs are exhausted and thin, fighting for survival, and have to spend extra effort to "reconstitute". The brumation of Gila monsters begins in October. Gila monsters can live up to 40 years in captivity, though rarely.
Little is known about the social behavior of Gila monster, but it has been observed engaging in male to male combat, in which the dominant male lies on top of the subordinate one and pins it with its front and hind limbs. While fighting, both lizards arch their bodies, pushing against each other and twisting around in an effort to gain the dominant position. A "wrestling match" ends when the pressure exerts their forces, although bouts may be repeated. These bouts are typically observed in the mating season. Males with greater strength and endurance are thought to enjoy greater reproductive success. Although the Gila monster has a low metabolism and one of the lowest lizard sprint speeds, it has one of the highest aerobic scope values (the increase in oxygen consumption from rest to maximum metabolic exertion) among lizards, allowing it to engage in intense aerobic activity for a sustained period of time.
Conservation status
Gila monsters are listed as near threatened by the IUCN. They are listed as "Apparently Secure" by NatureServe.
In 1952, the Gila monster became the first venomous animal to be given legal protection. They are now protected in all states of their distribution. International trade in the species is regulated under Appendix II of CITES.
Relocation
"Possibly the greatest threat to the continued existence of helodermatids is the man-made destruction of their habitat as the land is developed for construction or to create more cultivable land." Gila monsters found in these situations and relocated – with best intentions – up to away, return to where they were found within 2 months and at great effort. This is up to five times the normal energy use than if they had not been removed, which uses up their energy stores unnecessarily. The same is true for animals relocated to appropriate habitats. Besides this, they also become more exposed to predators. Therefore, the process of simple relocation is "naïve" and potentially dangerous for both the relocated animals and existing populations and for the inhabitants of the region where the resettlement is taking place. If relocating the lizards further away, they might be totally disoriented, thus their survival is still very questionable. A more successful strategy would be, for example, if the new "settlers" were offered intensive education about this species (e.g., limited toxicity, lifestyle) with the aim of tolerating the reptile or even being proud of having this unique "roommate" in one's own neighborhood.
In 1963, the San Diego Zoo became the first zoo to successfully breed Gila monsters in captivity. In the last two decades, experienced breeders have shared their knowledge and expertise to give advice to other herpetologists on overcoming the difficulties in Heloderma reproduction under human care.
Relationship with humans
Though the Gila monster is venomous, it poses little threat to humans due to its sluggish nature. Nevertheless, it has a fearsome reputation and is often killed by humans. Myths that have formed about the Gila monster include that the animal's breath is toxic enough to kill humans, that it can spit venom like a spitting cobra, that it can leap several feet in the air to attack, and that the Gila monster did not have an anus and therefore expelled waste from its mouth, the source of its venom and "fetid breath" (perhaps stemming from the fact that its venom in fact has an intense, specific smell). Among Native American tribes, the Gila monster has a varied reputation. The Apache believed its breath could kill a man, and the Tohono O'Odham and the Pima believed it possessed a spiritual power that could cause sickness. In contrast, the Seri and the Yaqui believed the Gila monster's hide had healing properties. The Navajo/Dine believe the Gila monster was the first medicine man.
In popular culture
The Gila monster starred as a monster in the film The Giant Gila Monster (though the titular monster was actually portrayed by a Mexican beaded lizard). It played a minor role in the motion picture The Treasure of the Sierra Madre. In Brock Brower's 1971 novel The Late Great Creature, fictional horror movie star Simon Moro is presented as famous for playing the reptilian werewolf-like Gila Man. The 2011 animated film Rango featured a Gila monster as an Old West outlaw named Bad Bill, voiced by Ray Winstone.
The Gila monster has also seen usage as a mascot and state symbol. The official mascot of Eastern Arizona College located in Thatcher, Arizona, is Gila Hank, a gun-toting, cowboy hat-wearing Gila monster. In 2017, the Vegas Golden Knights selected a Gila monster named Chance as their official mascot. In 2019, the state of Utah made the Gila monster its official state reptile.
In 2023, Australian band King Gizzard and the Lizard Wizard released a single titled "Gila Monster" from their album PetroDragonic Apocalypse, with the album's artwork featuring a Gila monster on its cover art and heavily within the album's concept narrative.
Gallery
| Biology and health sciences | Reptiles | null |
18953024 | https://en.wikipedia.org/wiki/Paleocene | Paleocene | The Paleocene ( ), or Palaeocene, is a geological epoch that lasted from about 66 to 56 million years ago (mya). It is the first epoch of the Paleogene Period in the modern Cenozoic Era. The name is a combination of the Ancient Greek palaiós meaning "old" and the Eocene Epoch (which succeeds the Paleocene), translating to "the old part of the Eocene".
The epoch is bracketed by two major events in Earth's history. The K–Pg extinction event, brought on by an asteroid impact (Chicxulub impact) and possibly volcanism (Deccan Traps), marked the beginning of the Paleocene and killed off 75% of species, most famously the non-avian dinosaurs. The end of the epoch was marked by the Paleocene–Eocene Thermal Maximum (PETM), which was a major climatic event wherein about 2,500–4,500 gigatons of carbon were released into the atmosphere and ocean systems, causing a spike in global temperatures and ocean acidification.
In the Paleocene, the continents of the Northern Hemisphere were still connected via some land bridges; and South America, Antarctica, and Australia had not completely separated yet. The Rocky Mountains were being uplifted, the Americas had not yet joined, the Indian Plate had begun its collision with Asia, and the North Atlantic Igneous Province was forming in the third-largest magmatic event of the last 150 million years. In the oceans, the thermohaline circulation probably was much different from what it is today, with downwellings occurring in the North Pacific rather than the North Atlantic, and water density mainly being controlled by salinity rather than temperature.
The K–Pg extinction event caused a floral and faunal turnover of species, with previously abundant species being replaced by previously uncommon ones. In the Paleocene, with a global average temperature of about , compared to in more recent times, the Earth had a greenhouse climate without permanent ice sheets at the poles, like the preceding Mesozoic. As such, there were forests worldwide—including at the poles—but they had low species richness in regards to plant life, and were populated by mainly small creatures that were rapidly evolving to take advantage of the recently emptied Earth. Though some animals attained great size, most remained rather small. The forests grew quite dense in the general absence of large herbivores. Mammals proliferated in the Paleocene, and the earliest placental and marsupial mammals are recorded from this time, but most Paleocene taxa have ambiguous affinities. In the seas, ray-finned fish rose to dominate open ocean and recovering reef ecosystems.
Etymology
The word "Paleocene" was first used by French paleobotanist and geologist Wilhelm Philipp Schimper in 1874 while describing deposits near Paris (spelled in his treatise). By this time, Italian geologist Giovanni Arduino had divided the history of life on Earth into the Primary (Paleozoic), Secondary (Mesozoic), and Tertiary in 1759; French geologist Jules Desnoyers had proposed splitting off the Quaternary from the Tertiary in 1829; and Scottish geologist Charles Lyell (ignoring the Quaternary) had divided the Tertiary Epoch into the Eocene, Miocene, Pliocene, and New Pliocene (Holocene) Periods in 1833. British geologist John Phillips had proposed the Cenozoic in 1840 in place of the Tertiary, and Austrian paleontologist Moritz Hörnes had introduced the Paleogene for the Eocene and Neogene for the Miocene and Pliocene in 1853. After decades of inconsistent usage, the newly formed International Commission on Stratigraphy (ICS), in 1969, standardized stratigraphy based on the prevailing opinions in Europe: the Cenozoic Era subdivided into the Tertiary and Quaternary sub-eras, and the Tertiary subdivided into the Paleogene and Neogene Periods. In 1978, the Paleogene was officially defined as the Paleocene, Eocene, and Oligocene Epochs; and the Neogene as the Miocene and Pliocene Epochs. In 1989, Tertiary and Quaternary were removed from the time scale due to the arbitrary nature of their boundary, but Quaternary was reinstated in 2009.
The term "Paleocene" is a portmanteau combination of the Ancient Greek palaios meaning "old", and the word "Eocene", and so means "the old part of the Eocene". The Eocene, in turn, is derived from Ancient Greek eo—eos meaning "dawn", and—cene kainos meaning "new" or "recent", as the epoch saw the dawn of recent, or modern, life. Paleocene did not come into broad usage until around 1920. In North America and mainland Europe, the standard spelling is "Paleocene", whereas it is "Palaeocene" in the UK. Geologist T. C. R. Pulvertaft has argued that the latter spelling is incorrect because this would imply either a translation of "old recent" or a derivation from "pala" and "Eocene", which would be incorrect because the prefix palæo- uses the ligature æ instead of "a" and "e" individually, so only both characters or neither should be dropped, not just one.
Geology
Boundaries
The Paleocene Epoch is the 10 million year time interval directly after the K–Pg extinction event, which ended the Cretaceous Period and the Mesozoic Era, and initiated the Cenozoic Era and the Paleogene Period. It is divided into three ages: the Danian spanning 66 to 61.6 million years ago (mya), the Selandian spanning 61.6 to 59.2 mya, and the Thanetian spanning 59.2 to 56 mya. It is succeeded by the Eocene.
The K–Pg boundary is clearly defined in the fossil record in numerous places around the world by a high-iridium band, as well as discontinuities with fossil flora and fauna. It is generally thought that a wide asteroid impact, forming the Chicxulub Crater in the Yucatán Peninsula in the Gulf of Mexico, and Deccan Trap volcanism caused a cataclysmic event at the boundary resulting in the extinction of 75% of all species.
The Paleocene ended with the Paleocene–Eocene thermal maximum, a short period of intense warming and ocean acidification brought about by the release of carbon en masse into the atmosphere and ocean systems, which led to a mass extinction of 30–50% of benthic foraminifera–planktonic species which are used as bioindicators of the health of a marine ecosystem—one of the largest in the Cenozoic. This event happened around 55.8 mya, and was one of the most significant periods of global change during the Cenozoic.
Stratigraphy
Geologists divide the rocks of the Paleocene into a stratigraphic set of smaller rock units called stages, each formed during corresponding time intervals called ages. Stages can be defined globally or regionally. For global stratigraphic correlation, the ICS ratify global stages based on a Global Boundary Stratotype Section and Point (GSSP) from a single formation (a stratotype) identifying the lower boundary of the stage. In 1989, the ICS decided to split the Paleocene into three stages: the Danian, Selandian, and Thanetian.
The Danian was first defined in 1847 by German-Swiss geologist Pierre Jean Édouard Desor based on the Danish chalks at Stevns Klint and Faxse, and was part of the Cretaceous, succeeded by the Tertiary Montian Stage. In 1982, after it was shown that the Danian and the Montian are the same, the ICS decided to define the Danian as starting with the K–Pg boundary, thus ending the practice of including the Danian in the Cretaceous. In 1991, the GSSP was defined as a well-preserved section in the El Haria Formation near El Kef, Tunisia, , and the proposal was officially published in 2006.
The Selandian and Thanetian are both defined in Itzurun beach by the Basque town of Zumaia, , as the area is a continuous early Santonian to early Eocene sea cliff outcrop. The Paleocene section is an essentially complete, exposed record thick, mainly composed of alternating hemipelagic sediments deposited at a depth of about . The Danian deposits are sequestered into the Aitzgorri Limestone Formation, and the Selandian and early Thanetian into the Itzurun Formation. The Itzurun Formation is divided into groups A and B corresponding to the two stages respectively. The two stages were ratified in 2008, and this area was chosen because of its completion, low risk of erosion, proximity to the original areas the stages were defined, accessibility, and the protected status of the area due to its geological significance.
The Selandian was first proposed by Danish geologist Alfred Rosenkrantz in 1924 based on a section of fossil-rich glauconitic marls overlain by gray clay which unconformably overlies Danian chalk and limestone. The area is now subdivided into the Æbelø Formation, Holmehus Formation, and the Østerrende Clay. The beginning of this stage was defined by the end of carbonate rock deposition from an open ocean environment in the North Sea region (which had been going on for the previous 40 million years). The Selandian deposits in this area are directly overlain by the Eocene Fur Formation—the Thanetian was not represented here—and this discontinuity in the deposition record is why the GSSP was moved to Zumaia. Today, the beginning of the Selandian is marked by the appearances of the nannofossils Fasciculithus tympaniformis, Neochiastozygus perfectus, and Chiasmolithus edentulus, though some foraminifera are used by various authors.
The Thanetian was first proposed by Swiss geologist Eugène Renevier, in 1873; he included the south England Thanet, Woolwich, and Reading formations. In 1880, French geologist Gustave Frédéric Dollfus narrowed the definition to just the Thanet Formation. The Thanetian begins a little after the mid-Paleocene biotic event—a short-lived climatic event caused by an increase in methane—recorded at Itzurun as a dark interval from a reduction of calcium carbonate. At Itzurun, it begins about above the base of the Selandian, and is marked by the first appearance of the algae Discoaster and a diversification of Heliolithus, though the best correlation is in terms of paleomagnetism. A chron is the occurrence of a geomagnetic reversal—when the North and South poles switch polarities. Chron 1 (C1n) is defined as modern day to about 780,000 years ago, and the n denotes "normal" as in the polarity of today, and an r "reverse" for the opposite polarity. The beginning of the Thanetian is best correlated with the C26r/C26n reversal.
Mineral and hydrocarbon deposits
Several economically important coal deposits formed during the Paleocene, such as the sub-bituminous Fort Union Formation in the Powder River Basin of Wyoming and Montana, which produces 43% of American coal; the Wilcox Group in Texas, the richest deposits of the Gulf Coastal Plain; and the Cerrejón mine in Colombia, the largest open-pit mine in the country. Paleocene coal has been mined extensively in Svalbard, Norway, since near the beginning of the 20th century, and late Paleocene and early Eocene coal is widely distributed across the Canadian Arctic Archipelago and northern Siberia. In the North Sea, Paleocene-derived natural gas reserves, when they were discovered, totaled approximately 2.23 trillion m3 (7.89 trillion ft3), and oil in place 13.54 billion barrels. Important phosphate deposits—predominantly of francolite—near Métlaoui, Tunisia were formed from the late Paleocene to the early Eocene.
Impact craters
Impact craters formed in the Paleocene include: the Connolly Basin crater in Western Australia less than 60 mya, the Texan Marquez crater 58 mya, the Greenlandic Hiawatha Glacier crater 58 mya, and possibly the Jordan Jabel Waqf as Suwwan crater which dates to between 56 and 37 mya. Vanadium-rich osbornite from the Isle of Skye, Scotland, dating to 60 mya may be impact ejecta. Craters were also formed near the K–Pg boundary, the largest the Mexican Chicxulub crater whose impact was a major precipitator of the K–Pg extinction, and also the Ukrainian Boltysh crater, dated to 65.4 mya the Canadian Eagle Butte crater (though it may be younger), the Vista Alegre crater (though this may date to about 115 mya). Silicate glass spherules along the Atlantic coast of the U.S. indicate a meteor impact in the region at the PETM.
Paleogeography
Paleotectonics
During the Paleocene, the continents continued to drift toward their present positions. In the Northern Hemisphere, the former components of Laurasia (North America and Eurasia) were, at times, connected via land bridges: Beringia (at 65.5 and 58 mya) between North America and East Asia, the De Geer route (from 71 to 63 mya) between Greenland and Scandinavia, the Thulean route (at 57 and 55.8 mya) between North America and Western Europe via Greenland, and the Turgai route connecting Europe with Asia (which were otherwise separated by the Turgai Strait at this time).
The Laramide orogeny, which began in the Late Cretaceous, continued to uplift the Rocky Mountains; it ended at the end of the Paleocene. Because of this and a drop in sea levels resulting from tectonic activity, the Western Interior Seaway, which had divided the continent of North America for much of the Cretaceous, had receded.
Between about 60.5 and 54.5 mya, there was heightened volcanic activity in the North Atlantic region—the third largest magmatic event in the last 150 million years—creating the North Atlantic Igneous Province. The proto-Iceland hotspot is sometimes cited as being responsible for the initial volcanism, though rifting and resulting volcanism have also contributed. This volcanism may have contributed to the opening of the North Atlantic Ocean and seafloor spreading, the divergence of the Greenland Plate from the North American Plate, and, climatically, the PETM by dissociating methane clathrate crystals on the seafloor resulting in the mass release of carbon.
North and South America remained separated by the Central American Seaway, though an island arc (the South Central American Arc) had already formed about 73 mya. The Caribbean Large Igneous Province (now the Caribbean Plate), which had formed from the Galápagos hotspot in the Pacific in the latest Cretaceous, was moving eastward as the North American and South American plates were getting pushed in the opposite direction due to the opening of the Atlantic (strike-slip tectonics). This motion would eventually uplift the Isthmus of Panama by 2.6 mya. The Caribbean Plate continued moving until about 50 mya when it reached its current position.
The components of the former southern supercontinent Gondwanaland in the Southern Hemisphere continued to drift apart, but Antarctica was still connected to South America and Australia. Africa was heading north towards Europe, and the Indian subcontinent towards Asia, which would eventually close the Tethys Ocean. The Indian and Eurasian Plates began colliding in the Paleocene, with uplift (and a land connection) beginning in the Miocene about 24–17 mya. There is evidence that some plants and animals could migrate between India and Asia during the Paleocene, possibly via intermediary island arcs.
Paleoceanography
In the modern thermohaline circulation, warm tropical water becomes colder and saltier at the poles and sinks (downwelling or deep water formation) that occurs at the North Atlantic near the North Pole and the Southern Ocean near the Antarctic Peninsula. In the Paleocene, the waterways between the Arctic Ocean and the North Atlantic were somewhat restricted, so North Atlantic Deep Water (NADW) and the Atlantic Meridional Overturning Circulation (AMOC)—which circulates cold water from the Arctic towards the equator—had not yet formed, and so deep water formation probably did not occur in the North Atlantic. The Arctic and Atlantic would not be connected by sufficiently deep waters until the early to middle Eocene.
There is evidence of deep water formation in the North Pacific to at least a depth of about . The elevated global deep water temperatures in the Paleocene may have been too warm for thermohaline circulation to be predominately heat driven. It is possible that the greenhouse climate shifted precipitation patterns, such that the Southern Hemisphere was wetter than the Northern, or the Southern experienced less evaporation than the Northern. In either case, this would have made the Northern more saline than the Southern, creating a density difference and a downwelling in the North Pacific traveling southward. Deep water formation may have also occurred in the South Atlantic.
It is largely unknown how global currents could have affected global temperature. The formation of the Northern Component Waters by Greenland in the Eocene—the predecessor of the AMOC—may have caused an intense warming in the North Hemisphere and cooling in the Southern, as well as an increase in deep water temperatures. In the PETM, it is possible deep water formation occurred in saltier tropical waters and moved polewards, which would increase global surface temperatures by warming the poles. Also, Antarctica was still connected to South America and Australia, and, because of this, the Antarctic Circumpolar Current—which traps cold water around the continent and prevents warm equatorial water from entering—had not yet formed. Its formation may have been related in the freezing of the continent. Warm coastal upwellings at the poles would have inhibited permanent ice cover. Conversely, it is possible deep water circulation was not a major contributor to the greenhouse climate, and deep water temperatures more likely change as a response to global temperature change rather than affecting it.
In the Arctic, coastal upwelling may have been largely temperature and wind-driven. In summer, the land surface temperature was probably higher than oceanic temperature, and the opposite was true in the winter, which is consistent with monsoon seasons in Asia. Open-ocean upwelling may have also been possible.
Climate
Average climate
The Paleocene climate was, much like in the Cretaceous, tropical or subtropical, and the poles were temperate, with an average global temperature of roughly . For comparison, the average global temperature for the period between 1951 and 1980 was . The latitudinal temperature gradient was approximately 0.24 °C per degree of latitude. The poles also lacked ice caps, though some alpine glaciation did occur in the Transantarctic Mountains.
The poles probably had a cool temperate climate; northern Antarctica, Australia, the southern tip of South America, what is now the US and Canada, eastern Siberia, and Europe warm temperate; middle South America, southern and northern Africa, South India, Middle America, and China arid; and northern South America, central Africa, North India, middle Siberia, and what is now the Mediterranean Sea tropical. South-central North America had a humid, monsoonal climate along its coastal plain, but conditions were drier to the west and at higher altitudes. Svalbard was temperate, having a mean temperature of 19.2 ± 2.49 °C during its warmest month and 1.7 ± 3.24 °С during its coldest.
Global deep water temperatures in the Paleocene likely ranged from , compared to in modern day. Based on the upper limit, average sea surface temperatures (SSTs) at 60°N and S would have been the same as deep sea temperatures, at 30°N and S about , and at the equator about . In the Danish Palaeocene sea, SSTs were cooler than those of the preceding Late Cretaceous and the succeeding Eocene. The Paleocene foraminifera assemblage globally indicates a defined deep-water thermocline (a warmer mass of water closer to the surface sitting on top of a colder mass nearer the bottom) persisting throughout the epoch. The Atlantic foraminifera indicate a general warming of sea surface temperature–with tropical taxa present in higher latitude areas–until the Late Paleocene when the thermocline became steeper and tropical foraminifera retreated back to lower latitudes.
Early Paleocene atmospheric CO2 levels at what is now Castle Rock, Colorado, were calculated to be between 352 and 1,110 parts per million (ppm), with a median of 616 ppm. Based on this and estimated plant-gas exchange rates and global surface temperatures, the climate sensitivity was calculated to be +3 °C when CO2 levels doubled, compared to 7 °C following the formation of ice at the poles. CO2 levels alone may have been insufficient in maintaining the greenhouse climate, and some positive feedbacks must have been active, such as some combination of cloud, aerosol, or vegetation related processes. A 2019 study identified changes in orbital eccentricity as the dominant drivers of climate between the late Cretaceous and the early Eocene.
Climatic events
The effects of the meteor impact and volcanism 66 mya and the climate across the K–Pg boundary were likely fleeting, and climate reverted to normal in a short time frame. The freezing temperatures probably reversed after three years and returned to normal within decades, sulfuric acid aerosols causing acid rain probably dissipated after 10 years, and dust from the impact blocking out sunlight and inhibiting photosynthesis would have lasted up to a year though potential global wildfires raging for several years would have released more particulates into the atmosphere. For the following half million years, the carbon isotope gradient—a difference in the 13C/12C ratio between surface and deep ocean water, causing carbon to cycle into the deep sea—may have shut down. This, termed a "Strangelove ocean", indicates low oceanic productivity; resultant decreased phytoplankton activity may have led to a reduction in cloud seeds and, thus, marine cloud brightening, causing global temperatures to increase by 6 °C (CLAW hypothesis). Following the extreme disruptions in the aftermath of the K-Pg extinction event, the relatively cool, though still greenhouse, conditions of the Late Cretaceous–Early Palaeogene Cool Interval (LKEPCI) that began in the Late Cretaceous continued.
The Dan–C2 Event 65.2 mya in the early Danian spanned about 100,000 years, and was characterized by an increase in carbon, particularly in the deep sea. Since the mid-Maastrichtian, more and more carbon had been sequestered in the deep sea possibly due to a global cooling trend and increased circulation into the deep sea. The Dan–C2 event may represent a release of this carbon after deep sea temperatures rose to a certain threshold, as warmer water can dissolve less carbon. Alternatively, the cause of the Dan-C2 event may have been a pulse of Deccan Traps volcanism. Savanna may have temporarily displaced forestland in this interval.
Around 62.2 mya in the late Danian, there was a warming event and evidence of ocean acidification associated with an increase in carbon; at this time, there was major seafloor spreading in the Atlantic and volcanic activity along the southeast margin of Greenland. The Latest Danian Event, also known as the Top Chron C27n Event, lasted about 200,000 years and resulted in a 1.6–2.8 °C increase in temperatures throughout the water column. Though the temperature in the latest Danian varied at about the same magnitude, this event coincides with an increase of carbon.
About 60.5 mya at the Danian/Selandian boundary, there is evidence of anoxia spreading out into coastal waters, and a drop in sea levels which is most likely explained as an increase in temperature and evaporation, as there was no ice at the poles to lock up water.
During the mid-Palaeocene biotic event (MPBE), also known as the Early Late Palaeocene Event (ELPE), around 59 Ma (roughly 50,000 years before the Selandian/Thanetian boundary), the temperature spiked probably due to a mass release of the deep sea methane hydrate into the atmosphere and ocean systems. Carbon was probably output for 10–11,000 years, and the aftereffects likely subsided around 52–53,000 years later. There is also evidence this occurred again 300,000 years later in the early Thanetian dubbed MPBE-2. Respectively, about 83 and 132 gigatons of methane-derived carbon were ejected into the atmosphere, which suggests a rise in temperature, and likely caused heightened seasonality and less stable environmental conditions. It may have also caused an increase of grass in some areas.
From 59.7 to 58.1 Ma, during the late Selandian and early Thanetian, organic carbon burial resulted in a period of climatic cooling, sea level fall and transient ice growth. This interval saw the highest δ18O values of the epoch.
Paleocene–Eocene Thermal Maximum
The Paleocene–Eocene Thermal Maximum was an approximately 200,000-year-long event where the global average temperature rose by some , and mid-latitude and polar areas may have exceeded modern tropical temperatures of . This was due to an ejection of 2,500–4,500 gigatons of carbon into the atmosphere, most commonly explained as the perturbation and release of methane clathrate deposits in the North Atlantic from tectonic activity and resultant increase in bottom water temperatures. Other proposed hypotheses include methane release from the heating of organic matter at the seafloor rather than methane clathrates, or melting permafrost.
The duration of carbon output is controversial, but most likely about 2,500 years. This carbon also interfered with the carbon cycle and caused ocean acidification, and potentially altered and slowed down ocean currents, the latter leading to the expansion of oxygen minimum zones (OMZs) in the deep sea. In surface water, OMZs could have also been caused from the formation of strong thermoclines preventing oxygen inflow, and higher temperatures equated to higher productivity leading to higher oxygen usurpation. Further, expanding OMZs could have caused the proliferation of sulfate-reducing microorganisms which create highly toxic hydrogen sulfide H2S as a waste product. During the event, the volume of sulfidic water may have been 10–20% of total ocean volume, in comparison to today's 1%. This may have also caused chemocline upwellings along continents and the dispersal of H2S into the atmosphere. During the PETM there was a temporary dwarfing of mammals apparently caused by the upward excursion in temperature.
Flora
The warm, wet climate supported tropical and subtropical forests worldwide, mainly populated by conifers and broad-leafed trees. In Patagonia, the landscape supported tropical rainforests, cloud rainforests, mangrove forests, swamp forests, savannas, and sclerophyllous forests. In the Colombian Cerrejón Formation, fossil flora belong to the same families as modern day flora—such as palm trees, legumes, aroids, and malvales—and the same is true in the North Dakotan Almont/Beicegel Creek—such as Ochnaceae, Cyclocarya, and Ginkgo cranei—indicating the same floral families have characterized South American rainforests and the American Western Interior since the Paleocene.
The extinction of large herbivorous dinosaurs may have allowed the forests to grow quite dense, and there is little evidence of wide open plains. Plants evolved several techniques to cope with high plant density, such as buttressing to better absorb nutrients and compete with other plants, increased height to reach sunlight, larger diaspore in seeds to provide added nutrition on the dark forest floor, and epiphytism where a plant grows on another plant in response to less space on the forest floor. Despite increasing density—which could act as fuel—wildfires decreased in frequency from the Cretaceous to the early Eocene as the atmospheric oxygen levels decreased to modern day levels, though they may have been more intense.
Recovery
There was a major die-off of plant species over the boundary; for example, in the Williston Basin of North Dakota, an estimated 1/3 to 3/5 of plant species went extinct. The K–Pg extinction event ushered in a floral turnover; for example, the once commonplace Araucariaceae conifers were almost fully replaced by Podocarpaceae conifers, and the Cheirolepidiaceae, a group of conifers that had dominated during most of the Mesozoic but had become rare during the Late Cretaceous became dominant trees in Patagonia, before going extinct. Some plant communities, such as those in eastern North America, were already experiencing an extinction event in the late Maastrichtian, particularly in the 1 million years before the K–Pg extinction event. The "disaster plants" that refilled the emptied landscape crowded out many Cretaceous plants, and resultantly, many went extinct by the middle Paleocene.
The strata immediately overlaying the K–Pg extinction event are especially rich in fern fossils. Ferns are often the first species to colonize areas damaged by forest fires, so this "fern spike" may mark the recovery of the biosphere following the impact (which caused blazing fires worldwide). The diversifying herb flora of the early Paleocene either represent pioneer species which re-colonized the recently emptied landscape, or a response to the increased amount of shade provided in a forested landscape. Lycopods, ferns, and angiosperm shrubs may have been important components of the Paleocene understory.
In general, the forests of the Paleocene were species-poor, and diversity did not fully recover until the end of the Paleocene. For example, the floral diversity of what is now the Holarctic region (comprising most of the Northern Hemisphere) was mainly early members of Ginkgo, Metasequoia, Glyptostrobus, Macginitiea, Platanus, Carya, Ampelopsis, and Cercidiphyllum. Patterns in plant recovery varied significantly with latitude, climate, and altitude. For example, what is now Castle Rock, Colorado featured a rich rainforest only 1.4 million years after the event, probably due to a rain shadow effect causing regular monsoon seasons. Conversely, low plant diversity and a lack of specialization in insects in the Colombian Cerrejón Formation, dated to 58 mya, indicates the ecosystem was still recovering from the K–Pg extinction event 7 million years later.
Angiosperms
Flowering plants (angiosperms), which had become dominant among forest taxa by the middle Cretaceous 110–90 mya, continued to develop and proliferate, more so to take advantage of the recently emptied niches and an increase in rainfall. Along with them coevolved the insects that fed on these plants and pollinated them. Predation by insects was especially high during the PETM. Many fruit-bearing plants appeared in the Paleocene in particular, probably to take advantage of the newly evolving birds and mammals for seed dispersal.
In what is now the Gulf Coast, angiosperm diversity increased slowly in the early Paleocene, and more rapidly in the middle and late Paleocene. This may have been because the effects of the K–Pg extinction event were still to some extent felt in the early Paleocene, the early Paleocene may not have had as many open niches, early angiosperms may not have been able to evolve at such an accelerated rate as later angiosperms, low diversity equates to lower evolution rates, or there was not much angiosperm migration into the region in the early Paleocene. Over the K–Pg extinction event, angiosperms had a higher extinction rate than gymnosperms (which include conifers, cycads, and relatives) and pteridophytes (ferns, horsetails, and relatives); zoophilous angiosperms (those that relied on animals for pollination) had a higher rate than anemophilous angiosperms; and evergreen angiosperms had a higher rate than deciduous angiosperms as deciduous plants can become dormant in harsh conditions.
In the Gulf Coast, angiosperms experienced another extinction event during the PETM, which they recovered quickly from in the Eocene through immigration from the Caribbean and Europe. During this time, the climate became warmer and wetter, and it is possible that angiosperms evolved to become stenotopic by this time, able to inhabit a narrow range of temperature and moisture; or, since the dominant floral ecosystem was a highly integrated and complex closed-canopy rainforest by the middle Paleocene, the plant ecosystems were more vulnerable to climate change. There is some evidence that, in the Gulf Coast, there was an extinction event in the late Paleocene preceding the PETM, which may have been due to the aforementioned vulnerability of complex rainforests, and the ecosystem may have been disrupted by only a small change in climate.
Polar forests
The warm Paleocene climate, much like that of the Cretaceous, allowed for diverse polar forests. Whereas precipitation is a major factor in plant diversity nearer the equator, polar plants had to adapt to varying light availability (polar nights and midnight suns) and temperatures. Because of this, plants from both poles independently evolved some similar characteristics, such as broad leaves. Plant diversity at both poles increased throughout the Paleocene, especially at the end, in tandem with the increasing global temperature.
At the North Pole, woody angiosperms had become the dominant plants, a reversal from the Cretaceous where herbs proliferated. The Iceberg Bay Formation on Ellesmere Island, Nunavut (latitude 75–80° N) shows remains of a late Paleocene dawn redwood forest, the canopy reaching around , and a climate similar to the Pacific Northwest. On the Alaska North Slope, Metasequoia was the dominant conifer. Much of the diversity represented migrants from nearer the equator. Deciduousness was dominant, probably to conserve energy by retroactively shedding leaves and retaining some energy rather than having them die from frostbite. In south-central Alaska, the Chickaloon Formation preserves peat-forming swamps dominated by taxodiaceous conifers and clastic floodplains occupied by angiosperm–conifer forests.
At the South Pole, due to the increasing isolation of Antarctica, many plant taxa were endemic to the continent instead of migrating down. Patagonian flora may have originated in Antarctica. The climate was much cooler than in the Late Cretaceous, though frost probably was not common in at least coastal areas. East Antarctica was likely warm and humid. Because of this, evergreen forests could proliferate as, in the absence of frost and a low probability of leaves dying, it was more energy efficient to retain leaves than to regrow them every year. One possibility is that the interior of the continent favored deciduous trees, though prevailing continental climates may have produced winters warm enough to support evergreen forests. As in the Cretaceous, podocarpaceous conifers, Nothofagus, and Proteaceae angiosperms were common.
Fauna
In the K–Pg extinction event, every land animal over was wiped out, leaving open several niches at the beginning of the epoch.
Mammals
Mammals had first appeared in the Late Triassic, and remained small and nocturnal throughout the Mesozoic to avoid competition with dinosaurs (nocturnal bottleneck), though, by the Middle Jurassic, they had branched out into several habitats—such as subterranean, arboreal, and aquatic— and the largest known Mesozoic mammal, Repenomamus robustus reached about in length and in weight–comparable to the modern day Virginia opossum. Though some mammals could sporadically venture out in daytime (cathemerality) by roughly 10 million years before the K–Pg extinction event, they only became strictly diurnal (active in the daytime) sometime after.
In general, Paleocene mammals retained this small size until near the end of the epoch, and, consequently, early mammal bones are not well preserved in the fossil record, and most of what is known comes from fossil teeth. Multituberculates, a now-extinct rodent-like group not closely related to any modern mammal, were the most successful group of mammals in the Mesozoic, and they reached peak diversity in the early Paleocene. During this time, multituberculate taxa had a wide range of dental complexity, which correlates to a broader range in diet for the group as a whole. Multituberculates declined in the late Paleocene and went extinct at the end of the Eocene, possibly due to competition from newly evolving rodents.
Nonetheless, following the K–Pg extinction event, mammals very quickly diversified and filled the empty niches. Mammal richness during this epoch, in contrast to the present day, varied insignificantly with latitude. Modern mammals are subdivided into therians (modern members are placentals and marsupials) and monotremes. These three groups all originated in the Cretaceous. Paleocene marsupials include Peradectes, and monotremes Monotrematum. The epoch featured the rise of many crown placental groups—groups that have living members in modern day—such as the earliest afrotherian Ocepeia, xenarthran Utaetus, rodent Tribosphenomys and Paramys, the forerunners of primates the Plesiadapiformes, earliest carnivorans Ravenictis and Pristinictis, possible pangolins Palaeanodonta, possible forerunners of odd-toed ungulates Phenacodontidae, and eulipotyphlans Nyctitheriidae. Though therian mammals had probably already begun to diversify around 10 to 20 million years before the K–Pg extinction event, average mammal size increased greatly after the boundary, and a radiation into frugivory (fruit-eating) and omnivory began, namely with the newly evolving large herbivores such as the Taeniodonta, Tillodonta, Pantodonta, Polydolopimorphia, and the Dinocerata. Large carnivores include the wolf-like Mesonychia, such as Ankalagon and Sinonyx.
Though there was an explosive diversification, the affinities of most Paleocene mammals are unknown, and only primates, carnivorans, and rodents have unambiguous Paleocene origins, resulting in a 10 million year gap in the fossil record of other mammalian crown orders. The most species-rich order of Paleocene mammals is Condylarthra, which is a wastebasket taxon for miscellaneous bunodont hoofed mammals. Other ambiguous orders include the Leptictida, Cimolesta, and Creodonta. This uncertainty blurs the early evolution of placentals.
Birds
According to DNA studies, modern birds (Neornithes) rapidly diversified following the extinction of the other dinosaurs in the Paleocene, and nearly all modern bird lineages can trace their origins to this epoch with the exception of fowl and palaeognaths. This was one of the fastest diversifications of any group, probably fueled by the diversification of fruit-bearing trees and associated insects, and the modern bird groups had likely already diverged within four million years of the K–Pg extinction event. However, the fossil record of birds in the Paleocene is rather poor compared to other groups, limited globally to mainly waterbirds such as the early penguin Waimanu. The earliest arboreal crown group bird known is Tsidiiyazhi, a mousebird dating to around 62 mya. The fossil record also includes early owls such as the large Berruornis from France, and the smaller Ogygoptynx from the United States.
Almost all archaic birds (any bird outside Neornithes) went extinct during the K–Pg extinction event, although the archaic Qinornis is recorded in the Paleocene. Their extinction may have led to the proliferation of neornithine birds in the Paleocene, and the only known Cretaceous neornithine bird is the waterbird Vegavis, and possibly also the waterbird Teviornis.
In the Mesozoic, birds and pterosaurs exhibited size-related niche partitioning—no known Late Cretaceous flying bird had a wingspan greater than nor exceeded a weight of , whereas contemporary pterosaurs ranged from , probably to avoid competition. Their extinction allowed flying birds to attain greater size, such as pelagornithids and pelecaniformes. The Paleocene pelagornithid Protodontopteryx was quite small compared to later members, with a wingspan of about , comparable to a gull. On the archipelago-continent of Europe, the flightless bird Gastornis was the largest herbivore at tall for the largest species, possibly due to lack of competition from newly emerging large mammalian herbivores which were prevalent on the other continents. The carnivorous terror birds in South America have a contentious appearance in the Paleocene with Paleopsilopterus, though the first definitive appearance is in the Eocene.
Reptiles
It is generally believed all non-avian dinosaurs went extinct at the K–Pg extinction event 66 mya, though there are a couple of controversial claims of Paleocene dinosaurs which would indicate a gradual decline of dinosaurs. Contentious dates include remains from the Hell Creek Formation dated 40,000 years after the boundary, and a hadrosaur femur from the San Juan Basin dated to 64.5 mya, but such stray late forms may be zombie taxa that were washed out and moved to younger sediments.
In the wake of the K–Pg extinction event, 83% of lizard and snake (squamate) species went extinct, and the diversity did not fully recover until the end of the Paleocene. However, since the only major squamate lineages to disappear in the event were the mosasaurs and polyglyphanodontians (the latter making up 40% of Maastrichtian lizard diversity), and most major squamate groups had evolved by the Cretaceous, the event probably did not greatly affect squamate evolution, and newly evolving squamates did not seemingly branch out into new niches as mammals. That is, Cretaceous and Paleogene squamates filled the same niches. Nonetheless, there was a faunal turnover of squamates, and groups that were dominant by the Eocene were not as abundant in the Cretaceous, namely the anguids, iguanas, night lizards, pythons, colubrids, boas, and worm lizards. Only small squamates are known from the early Paleocene—the largest snake Helagras was in length—but the late Paleocene snake Titanoboa grew to over long, the longest snake ever recorded. Kawasphenodon peligrensis from the early Paleocene of South America represents the youngest record of Rhynchocephalia outside of New Zealand, where the only extant representative of the order, the tuatara, resides.
Freshwater crocodilians and choristoderans were among the aquatic reptiles to have survived the K–Pg extinction event, probably because freshwater environments were not as impacted as marine ones. One example of a Paleocene crocodile is Borealosuchus, which averaged in length at the Wannagan Creek site. Among crocodyliformes, the aquatic and terrestrial dyrosaurs and the fully terrestrial sebecids would also survive the K-Pg extinction event, and a late surviving member of Pholidosauridae is also known from the Danian of Morocco. Three choristoderans are known from the Paleocene: The gharial-like neochoristoderans Champsosaurus—the largest is the Paleocene C. gigas at , Simoedosaurus—the largest specimen measuring , and an indeterminate species of the lizard like non-neochoristoderan Lazarussuchus around 44 centimetres in length. The last known choristoderes, belonging to the genus Lazarussuchus, are known from the Miocene.
Turtles experienced a decline in the Campanian (Late Cretaceous) during a cooling event, and recovered during the PETM at the end of the Paleocene. Turtles were not greatly affected by the K–Pg extinction event, and around 80% of species survived. In Colombia, a 60 million year old turtle with a carapace, Carbonemys, was discovered.
Amphibians
There is little evidence amphibians were affected very much by the K–Pg extinction event, probably because the freshwater habitats they inhabited were not as greatly impacted as marine environments. In the Hell Creek Formation of eastern Montana, a 1990 study found no extinction in amphibian species across the boundary. The true toads evolved during the Paleocene. The final record of albanerpetontids from North America and outside of Europe and Anatolia, an unnamed species of Albanerpeton, is known from the Paleocene aged Paskapoo Formation in Canada.
Fish
The small pelagic fish population recovered rather quickly, and there was a low extinction rate for sharks and rays. Overall, only 12% of fish species went extinct. During the Cretaceous, fishes were not very abundant, probably due to heightened predation by or competition with ammonites and squid, although large predatory fish did exist, including ichthyodectids, pachycormids and pachyrhizodontids. Almost immediately following the K–Pg extinction event, ray-finned fish — today, representing nearly half of all vertebrate taxa – became much more numerous and increased in size, and rose to dominate the open-oceans. Acanthomorphs—a group of ray-finned fish which, today, represent a third of all vertebrate life—experienced a massive diversification following the K–Pg extinction event, dominating marine ecosystems by the end of the Paleocene, refilling vacant, open-ocean predatory niches as well as spreading out into recovering reef systems. In specific, percomorphs diversified faster than any other vertebrate group at the time, with the exception of birds; Cretaceous percomorphs varied very little in body plan, whereas, by the Eocene, percomorphs evolved into vastly varying creatures such as early scombrids (today, tuna, mackerels, and bonitos), barracudas, jacks, billfish, flatfishes, and aulostomoid (trumpetfish and cornetfish). However, the discovery of the Cretaceous cusk eel Pastorius shows that the body plans of at least some percomorphs were already highly variable, perhaps indicating an already diverse array of percomorph body plans before the Paleocene.
Conversely, sharks and rays appear to have been unable to exploit the vacant niches, and recovered the same pre-extinction abundance. There was a faunal turnover of sharks from mackerel sharks to ground sharks, as ground sharks are more suited to hunting the rapidly diversifying ray-finned fish whereas mackerel sharks target larger prey. The first megatoothed shark, Otodus obliquus—the ancestor of the giant megalodon—is recorded from the Paleocene.
Several Paleocene freshwater fish are recorded from North America, including bowfins, gars, arowanas, Gonorynchidae, common catfish, smelts, and pike.
Insects and arachnids
Insect recovery varied from place to place. For example, it may have taken until the PETM for insect diversity to recover in the western interior of North America, whereas Patagonian insect diversity had recovered by four million years after the K–Pg extinction event. In some areas, such as the Bighorn Basin in Wyoming, there is a dramatic increase in plant predation during the PETM, although this is probably not indicative of a diversification event in insects due to rising temperatures because plant predation decreases following the PETM. More likely, insects followed their host plant or plants which were expanding into mid-latitude regions during the PETM, and then retreated afterward.
The middle-to-late Paleocene French Menat Formation shows an abundance of beetles (making up 77.5% of the insect diversity)—especially weevils (50% of diversity), jewel beetles, leaf beetles, and reticulated beetles—as well as other true bugs—such as pond skaters—and cockroaches. To a lesser degree, there are also orthopterans, hymenopterans, butterflies, and flies, though planthoppers were more common than flies. Representing less than 1% of fossil remains dragonflies, caddisflies, mayflies, earwigs, mantises, net-winged insects, and possibly termites.
The Wyoming Hanna Formation is the only known Paleocene formation to produce sizable pieces of amber, as opposed to only small droplets. The amber was formed by a single or a closely related group of either taxodiaceaen or pine tree(s) which produced cones similar to those of dammaras. Only one insect, a thrips, has been identified.
There is a gap in the ant fossil record from 78 to 55 mya, except for the aneuretine Napakimyrma paskapooensis from the 62–56 million year old Canadian Paskapoo Formation. Given high abundance in the Eocene, two of the modern dominant ant subfamilies—Ponerinae and Myrmicinae—likely originated and greatly diversified in the Paleocene, acting as major hunters of arthropods, and probably competed with each other for food and nesting grounds in the dense angiosperm leaf litter. Myrmicines expanded their diets to seeds and formed trophobiotic symbiotic relationships with aphids, mealybugs, treehoppers, and other honeydew secreting insects which were also successful in angiosperm forests, allowing them to invade other biomes, such as the canopy or temperate environments, and achieve a worldwide distribution by the middle Eocene.
About 80% of the butterfly and moth (lepidopteran) fossil record occurs in the early Paleogene, specifically the late Paleocene and the middle-to-late Eocene. Most Paleocene lepidopteran compression fossils come from the Danish Fur Formation. Though there is low family-level diversity in the Paleocene compared to later epochs, this may be due to a largely incomplete fossil record. The evolution of bats had a profound effect on lepidopterans, which feature several anti-predator adaptations such as echolocation jamming and the ability to detect bat signals.
Bees were likely heavily impacted by the K–Pg extinction event and a die-off of flowering plants, though the bee fossil record is very limited. The oldest kleptoparasitic bee, Paleoepeolus, is known from the Paleocene 60 mya.
Though the Eocene features, by far, the highest proportion of known fossil spider species, the Paleocene spider assemblage is quite low. Some spider groups began to diversify around the PETM, such as jumping spiders, and possibly coelotine spiders (members of the funnel weaver family).
The diversification of mammals had a profound effect on parasitic insects, namely the evolution of bats, which have more ectoparasites than any other known mammal or bird. The PETM's effect on mammals greatly impacted the evolution of fleas, ticks, and oestroids.
Marine invertebrates
Among marine invertebrates, plankton and those with a planktonic stage in their development (meroplankton) were most impacted by the K–Pg extinction event, and plankton populations crashed. Nearly 90% of all calcifying plankton species perished. This reverberated up and caused a global marine food chain collapse, namely with the extinction of ammonites and large raptorial marine reptiles. Nonetheless, the rapid diversification of large fish species indicates a healthy plankton population through the Paleocene.
Marine invertebrate diversity may have taken about 7 million years to recover, though this may be a preservation artifact as anything smaller than is unlikely to be fossilized, and body size may have simply decreased across the boundary. A 2019 study found that in Seymour Island, Antarctica, the marine life assemblage consisted primarily of burrowing creatures—such as burrowing clams and snails—for around 320,000 years after the K–Pg extinction event, and it took around a million years for the marine diversity to return to previous levels. Areas closer to the equator may have been more affected. Sand dollars first evolved in the late Paleocene. The Late Cretaceous decapod crustacean assemblage of James Ross Island appears to have been mainly pioneer species and the ancestors of modern fauna, such as the first Antarctic crabs and the first appearance of the lobsters of the genera Linuparus, Metanephrops, and Munidopsis which still inhabit Antarctica today.
In the Cretaceous, the main reef-building creatures were the box-like bivalve rudists instead of coral—though a diverse Cretaceous coral assemblage did exist—and rudists had collapsed by the time of the K–Pg extinction event. Some corals are known to have survived in higher latitudes in the Late Cretaceous and into the Paleogene, and hard coral-dominated reefs may have recovered by 8 million years after the K–Pg extinction event, though the coral fossil record of this time is rather sparse. Though there was a lack of extensive coral reefs in the Paleocene, there were some colonies—mainly dominated by zooxanthellate corals—in shallow coastal (neritic) areas. Starting in the latest Cretaceous and continuing until the early Eocene, calcareous corals rapidly diversified. Corals probably competed mainly with red and coralline algae for space on the seafloor. Calcified dasycladalean green algae experienced the greatest diversity in their evolutionary history in the Paleocene. Though coral reef ecosystems do not become particularly abundant in the fossil record until the Miocene (possibly due to preservation bias), strong Paleocene coral reefs have been identified in what are now the Pyrenees (emerging as early as 63 mya), with some smaller Paleocene coral reefs identified across the Mediterranean region.
| Physical sciences | Geological timescale | Earth science |
18953991 | https://en.wikipedia.org/wiki/Bagheera%20kiplingi | Bagheera kiplingi | Bagheera kiplingi is a species of jumping spider found in Central America, including Mexico, Costa Rica, and Guatemala. It is the type species of the genus Bagheera, which includes three other species, including B. prosper. B. kiplingi is notable for its peculiar diet, which is mostly herbivorous. No other known species of omnivorous spider has such a markedly herbivorous diet.
Taxonomy
The genus name is derived from Bagheera, the black panther from Rudyard Kipling's The Jungle Book, with the species name honoring Kipling himself. Other salticid genera with names of Kipling's characters are Akela, Messua, and Nagaina. All four were named by George and Elizabeth Peckham in 1896.
Only the male was described in 1896; the female was first described 100 years later by Wayne Maddison.
Description
Bagheera kiplingi is a colorful, sexually dimorphic species. The male has amber legs, a dark cephalothorax that is greenish in the upper region near the front, and a slender reddish abdomen with green transversal lines. The female's amber front legs are sturdier than the other, slender legs, which are light yellow. It has a reddish-brown cephalothorax with the top region near the front black. The female's rather large abdomen is light brown with dark brown and greenish markings.
Diet
Bagheera kiplingi inhabits Mimosoideae trees, Vachellia in particular, where it consumes specialized protein- and fat-rich nubs called Beltian bodies. The nubs form at the leaf tips of the acacia as part of a symbiotic relationship with certain species of ants. The spiders actively avoid the ants that attempt to guard the Beltian bodies (their food source) against intruders. Although the Beltian bodies account for over 90% of B. kiplingi diet, the spiders also consume nectar and occasionally steal ant larvae from passing worker ants for food. Sometimes, they cannibalize conspecifics, especially during the dry season.
Despite their occasional acts of predation, the spiders' tissues have been found to exhibit isotopic signatures typical of herbivorous animals, implying that most of their food comes from plants. The mechanism by which they process, ingest, and metabolize the Beltian bodies is still unresearched. The vast majority of spiders liquefy their prey using digestive enzymes before sucking it in.
The degree of herbivory varies depending on environment. In Mexico, B. kiplingi inhabit more than 50% of Vachellia collinsii trees and feed almost exclusively on an herbivorous diet. In Costa Rica, the B. kiplingi population inhabit less than 5% of Acacia trees and their diet is less herbivorous. Although this species is mostly territorial and forages solitarily, populations of several hundred specimens have been found on individual acacias in Mexico, with more than twice as many females as males. B. kiplingi appears to breed throughout the year. Observations of adult females guarding hatchlings and clutches suggest that the species is quasisocial.
| Biology and health sciences | Spiders | Animals |
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