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34713249 | https://en.wikipedia.org/wiki/Forensic%20geophysics | Forensic geophysics | Forensic geophysics is a branch of forensic science and is the study, the search, the localization and the mapping of buried objects or elements beneath the soil or the water, using geophysics tools for legal purposes. There are various geophysical techniques for forensic investigations in which the targets are buried and have different dimensions (from weapons or metallic barrels to human burials and bunkers). Geophysical methods have the potential to aid the search and the recovery of these targets because they can non-destructively and rapidly investigate large areas where a suspect, illegal burial or, in general, a forensic target is hidden in the subsoil. When in the subsurface there is a contrast of physical properties between a target and the material in which it is buried, it is possible to individuate and define precisely the concealing place of the searched target. It is also possible to recognize evidences of human soil occupation or excavation, both recent and older. Forensic geophysics is an evolving technique that is gaining popularity and prestige in law enforcement.
Searched for objects obviously include clandestine graves of murder victims, but also include unmarked burials in graveyards and cemeteries, weapons used in criminal activities and environmental crime illegally dumping material.
There are various near-surface geophysical techniques that can be utilised to detect a near-surface buried object, which should be site and case-specific. A thorough desk study (including historical maps), utility survey, site reconnaissance and control studies should be undertaken before trial geophysical surveys and then full geophysical surveys are undertaken in phased investigations. Note also other search techniques should be used to first to prioritise suspect areas, for example cadaver dogs or forensic geomorphologists.
Techniques
For large-scale buried objects, seismic surveys may be appropriate but these have, at best, 2m vertical resolution so may not be ideal for certain targets, more typically they are used to detect bedrock below the surface.
For relatively quick site surveys, bulk ground electrical conductivity surveys can be collected which identifies areas of disturbance of different ground but these can suffer from a lack of resolution. This recent Black Death investigation in central London shows an example. shows a successful woodland search for a cold case in woodland in New Zealand.
Ground-penetrating radar (or GPR) has a typical maximum depth below ground level (bgl) of 10 m, depending upon the antennae frequencies used, typically 50 MHz to 1.2 Gz. The higher the frequency the smaller the object that can be resolved but also penetration depths decrease, so operators need to think carefully when choosing antennae frequencies and, ideally, undertake trial surveys using different antennae over a target at a known depth onsite. GPR is the most popularly used technique in forensic search, but is not suitable in certain soil types and environments, e.g. coastal (i.e. salt-rich) and clay-rich soils (lack of penetration). 2D profiles can be relatively quickly collected and, if time permits, successive profiles can be used to generate 3D datasets which may resolve more subtle targets. Recent studies have used GPR to locate mass graves from the Spanish Civil War in mountainous and urban environments.
Electrical resistivity methods can also detect objects, especially in clay-rich soil which would preclude the use of GPR. There are different equipment configurations, the dipole-dipole (fixed-offset) method is the most common which can traverse across an area, measuring resistivity variations at a set depth (typically 1-2x probe separations) which have been used in forensic searches. More slower methods are putting out many probes and collecting both spatially horizontally and vertically, called Electrical resistivity imaging (ERI). Multiple 2D profiles is termed electrical resistivity tomography (ERT).
Magnetometry can detect buried metal (or indeed fired objects such as bricks or even where surface fires were) using simple total field magnetometers, through to fluxgate gradiometers and high-end alkali vapour gradiometers, depending upon accuracy (and cost) required. Surface magnetic susceptibility has also shown recent promise for forensic search.
Water-based searches are also becoming more common, with specialist marine magnetometers, side-scan sonar and other acoustic methods and even water-penetrating radar methods used to rapidly scan bottoms of ponds, lakes, rivers and near-shore depositional environments.
Controlled research
There has been recent efforts to undertake research over known buried and below-water surface simulated forensic targets in order to gain an insight into optimum search technique(s) and/or equipment configuration(s). Most commonly, this involved the burial porcine cadavers and long-term monitoring for soilwater, seasonal effects on electrical resistivity surveys, burial in walls and beneath concrete, and Long-Term monitoring in the UK, the US and Latin America. Finally there has been surveys in graveyards over graves of known ages to determine the geophysical responses of multi-geophysical techniques with increasing burial ages
| Physical sciences | Geophysics | Earth science |
2922492 | https://en.wikipedia.org/wiki/Ponte%20Sant%27Angelo | Ponte Sant'Angelo | Ponte Sant'Angelo, originally the Aelian Bridge or Pons Aelius, is a Roman bridge in Rome, Italy, completed in 134 AD by Roman Emperor Hadrian (Publius Aelius Hadrianus), to span the Tiber from the city centre to his newly constructed mausoleum, now the towering Castel Sant'Angelo. The bridge is faced with travertine marble and spans the Tiber with five arches, three of which are Roman; it was approached by means of a ramp from the river. The bridge is now solely pedestrian and provides a scenic view of Castel Sant'Angelo. It links the rioni of Ponte (which was named after the bridge itself), and Borgo, to which the bridge administratively belongs.
History
Starting with the early Middle Ages, the original name was forgotten: after the ruin of Nero's Bridge, pilgrims were forced to use this bridge to reach St Peter's Basilica, hence it was known also with the name of "bridge of Saint Peter" (pons Sancti Petri). In the sixth century, under Pope Gregory I, both the castle and the bridge took on the name Sant'Angelo, explained by a legend that an angel appeared on the roof of the castle to announce the end of the plague. Dante writes in his Divine Comedy that during the jubilee of 1300, due to the large number of pilgrims going and coming from Saint Peter, two separate lanes were arranged on the bridge. During the 1450 jubilee, balustrades of the bridge collapsed, due to the great crowds of the pilgrims, and many drowned in the river. In response, some houses at the head of the bridge as well as a Roman triumphal arch were pulled down in order to widen the route for pilgrims.
In 1535, Pope Clement VII allocated the toll income of the bridge to erecting the statues of the apostles Saint Peter (holding a book, with the pedestal inscription Rione XIV) by Lorenzetto, and Saint Paul (holding a broken sword and a book, with the pedestal inscription Borgo) by Paolo Romano to which subsequently the four evangelists and the patriarchs were added to other statues representing Adam, Noah, Abraham, and Moses. For centuries after the 16th century, the bridge was used to expose the bodies of those executed in the nearby Piazza di Ponte, at the left bridge head. In 1669 Pope Clement IX commissioned replacements for the aging stucco angels by Raffaello da Montelupo, commissioned by Paul III. Bernini's program, one of his last large projects, called for ten angels holding instruments of the Passion: he personally only finished the two originals of the Angel with the Superscription "I.N.R.I." and the Angel with the Crown of Thorns, but these were kept by Clement IX for his own pleasure. They are now in the church of Sant'Andrea delle Fratte, also in Rome.
At the end of the 19th century, due to the works for the construction of the Lungotevere, the two Roman ramps which linked the bridge with the two banks were destroyed, and in their place two arches similar to the Roman ones were built.
For the Great Jubilee in 2000, the Lungotevere on the right bank between the bridge and the castle became a pedestrian area.
On 21 October 2019 Austrian activist Alexander Tschugguel threw five Pachamama statues, which he had stolen from a display at Santa Maria in Traspontina as part of the Amazon Synod, off of the bridge into the Tiber.
List of angels
| Technology | Bridges | null |
2924436 | https://en.wikipedia.org/wiki/Electromagnetic%20wave%20equation | Electromagnetic wave equation | The electromagnetic wave equation is a second-order partial differential equation that describes the propagation of electromagnetic waves through a medium or in a vacuum. It is a three-dimensional form of the wave equation. The homogeneous form of the equation, written in terms of either the electric field or the magnetic field , takes the form:
where
is the speed of light (i.e. phase velocity) in a medium with permeability , and permittivity , and is the Laplace operator. In a vacuum, , a fundamental physical constant. The electromagnetic wave equation derives from Maxwell's equations. In most older literature, is called the magnetic flux density or magnetic induction. The following equationspredicate that any electromagnetic wave must be a transverse wave, where the electric field and the magnetic field are both perpendicular to the direction of wave propagation.
The origin of the electromagnetic wave equation
In his 1865 paper titled A Dynamical Theory of the Electromagnetic Field, James Clerk Maxwell utilized the correction to Ampère's circuital law that he had made in part III of his 1861 paper On Physical Lines of Force. In Part VI of his 1864 paper titled Electromagnetic Theory of Light, Maxwell combined displacement current with some of the other equations of electromagnetism and he obtained a wave equation with a speed equal to the speed of light. He commented:
The agreement of the results seems to show that light and magnetism are affections of the same substance, and that light is an electromagnetic disturbance propagated through the field according to electromagnetic laws.
Maxwell's derivation of the electromagnetic wave equation has been replaced in modern physics education by a much less cumbersome method involving combining the corrected version of Ampère's circuital law with Faraday's law of induction.
To obtain the electromagnetic wave equation in a vacuum using the modern method, we begin with the modern 'Heaviside' form of Maxwell's equations. In a vacuum- and charge-free space, these equations are:
These are the general Maxwell's equations specialized to the case with charge and current both set to zero.
Taking the curl of the curl equations gives:
We can use the vector identity
where is any vector function of space. And
where is a dyadic which when operated on by the divergence operator yields a vector. Since
then the first term on the right in the identity vanishes and we obtain the wave equations:
where
is the speed of light in free space.
Covariant form of the homogeneous wave equation
These relativistic equations can be written in contravariant form as
where the electromagnetic four-potential is
with the Lorenz gauge condition:
and where
is the d'Alembert operator.
Homogeneous wave equation in curved spacetime
The electromagnetic wave equation is modified in two ways, the derivative is replaced with the covariant derivative and a new term that depends on the curvature appears.
where is the Ricci curvature tensor and the semicolon indicates covariant differentiation.
The generalization of the Lorenz gauge condition in curved spacetime is assumed:
Inhomogeneous electromagnetic wave equation
Localized time-varying charge and current densities can act as sources of electromagnetic waves in a vacuum. Maxwell's equations can be written in the form of a wave equation with sources. The addition of sources to the wave equations makes the partial differential equations inhomogeneous.
Solutions to the homogeneous electromagnetic wave equation
The general solution to the electromagnetic wave equation is a linear superposition of waves of the form
for virtually well-behaved function of dimensionless argument , where is the angular frequency (in radians per second), and is the wave vector (in radians per meter).
Although the function can be and often is a monochromatic sine wave, it does not have to be sinusoidal, or even periodic. In practice, cannot have infinite periodicity because any real electromagnetic wave must always have a finite extent in time and space. As a result, and based on the theory of Fourier decomposition, a real wave must consist of the superposition of an infinite set of sinusoidal frequencies.
In addition, for a valid solution, the wave vector and the angular frequency are not independent; they must adhere to the dispersion relation:
where is the wavenumber and is the wavelength. The variable can only be used in this equation when the electromagnetic wave is in a vacuum.
Monochromatic, sinusoidal steady-state
The simplest set of solutions to the wave equation result from assuming sinusoidal waveforms of a single frequency in separable form:
where
is the imaginary unit,
is the angular frequency in radians per second,
is the frequency in hertz, and
is Euler's formula.
Plane wave solutions
Consider a plane defined by a unit normal vector
Then planar traveling wave solutions of the wave equations are
where is the position vector (in meters).
These solutions represent planar waves traveling in the direction of the normal vector . If we define the direction as the direction of , and the direction as the direction of , then by Faraday's Law the magnetic field lies in the direction and is related to the electric field by the relation
Because the divergence of the electric and magnetic fields are zero, there are no fields in the direction of propagation.
This solution is the linearly polarized solution of the wave equations. There are also circularly polarized solutions in which the fields rotate about the normal vector.
Spectral decomposition
Because of the linearity of Maxwell's equations in a vacuum, solutions can be decomposed into a superposition of sinusoids. This is the basis for the Fourier transform method for the solution of differential equations. The sinusoidal solution to the electromagnetic wave equation takes the form
where
is time (in seconds),
is the angular frequency (in radians per second),
is the wave vector (in radians per meter), and
is the phase angle (in radians).
The wave vector is related to the angular frequency by
where is the wavenumber and is the wavelength.
The electromagnetic spectrum is a plot of the field magnitudes (or energies) as a function of wavelength.
Multipole expansion
Assuming monochromatic fields varying in time as , if one uses Maxwell's Equations to eliminate , the electromagnetic wave equation reduces to the Helmholtz equation for :
with as given above. Alternatively, one can eliminate in favor of to obtain:
A generic electromagnetic field with frequency can be written as a sum of solutions to these two equations. The three-dimensional solutions of the Helmholtz Equation can be expressed as expansions in spherical harmonics with coefficients proportional to the spherical Bessel functions. However, applying this expansion to each vector component of or will give solutions that are not generically divergence-free (), and therefore require additional restrictions on the coefficients.
The multipole expansion circumvents this difficulty by expanding not or , but or into spherical harmonics. These expansions still solve the original Helmholtz equations for and because for a divergence-free field , . The resulting expressions for a generic electromagnetic field are:
where and are the electric multipole fields of order (l, m), and and are the corresponding magnetic multipole fields, and and are the coefficients of the expansion. The multipole fields are given by
where are the spherical Hankel functions, and are determined by boundary conditions, and
are vector spherical harmonics normalized so that
The multipole expansion of the electromagnetic field finds application in a number of problems involving spherical symmetry, for example antennae radiation patterns, or nuclear gamma decay. In these applications, one is often interested in the power radiated in the far-field. In this regions, the and fields asymptotically approach
The angular distribution of the time-averaged radiated power is then given by
| Physical sciences | Electromagnetic radiation | Physics |
2925371 | https://en.wikipedia.org/wiki/Radiative%20transfer | Radiative transfer | Radiative transfer (also called radiation transport) is the physical phenomenon of energy transfer in the form of electromagnetic radiation. The propagation of radiation through a medium is affected by absorption, emission, and scattering processes. The equation of radiative transfer describes these interactions mathematically. Equations of radiative transfer have application in a wide variety of subjects including optics, astrophysics, atmospheric science, and remote sensing. Analytic solutions to the radiative transfer equation (RTE) exist for simple cases but for more realistic media, with complex multiple scattering effects, numerical methods are required.
The present article is largely focused on the condition of radiative equilibrium.
Definitions
The fundamental quantity that describes a field of radiation is called spectral radiance in radiometric terms (in other fields it is often called specific intensity). For a very small area element in the radiation field, there can be electromagnetic radiation passing in both senses in every spatial direction through it. In radiometric terms, the passage can be completely characterized by the amount of energy radiated in each of the two senses in each spatial direction, per unit time, per unit area of surface of sourcing passage, per unit solid angle of reception at a distance, per unit wavelength interval being considered (polarization will be ignored for the moment).
In terms of the spectral radiance, , the energy flowing across an area element of area located at in time in the solid angle about the direction in the frequency interval to is
where is the angle that the unit direction vector makes with a normal to the area element. The units of the spectral radiance are seen to be energy/time/area/solid angle/frequency. In MKS units this would be W·m−2·sr−1·Hz−1 (watts per square-metre-steradian-hertz).
The equation of radiative transfer
The equation of radiative transfer simply says that as a beam of radiation travels, it loses energy to absorption, gains energy by emission processes, and redistributes energy by scattering. The differential form of the equation for radiative transfer is:
where is the speed of light, is the emission coefficient, is the scattering opacity, is the absorption opacity, is the mass density and the term represents radiation scattered from other directions onto a surface.
Solutions to the equation of radiative transfer
Solutions to the equation of radiative transfer form an enormous body of work. The differences however, are essentially due to the various forms for the emission and absorption coefficients. If scattering is ignored, then a general steady state solution in terms of the emission and absorption coefficients may be written:
where is the optical depth of the medium between positions and :
Local thermodynamic equilibrium
A particularly useful simplification of the equation of radiative transfer occurs under the conditions of local thermodynamic equilibrium (LTE). It is important to note that local equilibrium may apply only to a certain subset of particles in the system. For example, LTE is usually applied only to massive particles. In a radiating gas, the photons being emitted and absorbed by the gas do not need to be in a thermodynamic equilibrium with each other or with the massive particles of the gas in order for LTE to exist.
In this situation, the absorbing/emitting medium consists of massive particles which are locally in equilibrium with each other, and therefore have a definable temperature (Zeroth Law of Thermodynamics). The radiation field is not, however in equilibrium and is being entirely driven by the presence of the massive particles. For a medium in LTE, the emission coefficient and absorption coefficient are functions of temperature and density only, and are related by:
where is the black body spectral radiance at temperature T. The solution to the equation of radiative transfer is then:
Knowing the temperature profile and the density profile of the medium is sufficient to calculate a solution to the equation of radiative transfer.
The Eddington approximation
The Eddington approximation is distinct from the two-stream approximation. The two-stream approximation assumes that the intensity is constant with angle in the upward hemisphere, with a different constant value in the downward hemisphere. The Eddington approximation instead assumes that the intensity is a linear function of , i.e.
where is the normal direction to the slab-like medium. Note that expressing angular integrals in terms of simplifies things because appears in the Jacobian of integrals in spherical coordinates. The Eddington approximation can be used to obtain the spectral radiance in a "plane-parallel" medium (one in which properties only vary in the perpendicular direction) with isotropic frequency-independent scattering.
Extracting the first few moments of the spectral radiance with respect to yields
Thus the Eddington approximation is equivalent to setting . Higher order versions of the Eddington approximation also exist, and consist of more complicated linear relations of the intensity moments. This extra equation can be used as a closure relation for the truncated system of moments.
Note that the first two moments have simple physical meanings. is the isotropic intensity at a point, and is the flux through that point in the direction.
The radiative transfer through an isotropically scattering medium with scattering coefficient at local thermodynamic equilibrium is given by
Integrating over all angles yields
Premultiplying by , and then integrating over all angles gives
Substituting in the closure relation, and differentiating with respect to allows the two above equations to be combined to form the radiative diffusion equation
This equation shows how the effective optical depth in scattering-dominated systems may be significantly different from that given by the scattering opacity if the absorptive opacity is small.
| Physical sciences | Electromagnetic radiation | Physics |
2925388 | https://en.wikipedia.org/wiki/Temnospondyli | Temnospondyli | Temnospondyli (from Greek τέμνειν, temnein 'to cut' and σπόνδυλος, spondylos 'vertebra') or temnospondyls is a diverse ancient order of small to giant tetrapods—often considered primitive amphibians—that flourished worldwide during the Carboniferous, Permian and Triassic periods, with fossils being found on every continent. A few species continued into the Jurassic and Early Cretaceous periods, but all had gone extinct by the Late Cretaceous. During about 210 million years of evolutionary history, they adapted to a wide range of habitats, including freshwater, terrestrial, and even coastal marine environments. Their life history is well understood, with fossils known from the larval stage, metamorphosis and maturity. Most temnospondyls were semiaquatic, although some were almost fully terrestrial, returning to the water only to breed. These temnospondyls were some of the first vertebrates fully adapted to life on land. Although temnospondyls are amphibians, many had characteristics such as scales and large armour-like bony plates (osteoderms) that generally distinguish them from the modern soft-bodied lissamphibians (frogs and toads, newts, salamanders and caecilians).
Temnospondyls have been known since the early 19th century, and were initially thought to be reptiles. They were described at various times as batrachians, stegocephalians and labyrinthodonts, although these names are now rarely used. Animals now grouped in Temnospondyli were spread out among several amphibian groups until the early 20th century, when they were found to belong to a distinct taxon based on the structure of their vertebrae. Temnospondyli means "cut vertebrae", as each vertebra is divided into several parts (intercentrum, paired pleurocentra, neural arch), although this occurs widely among other early tetrapods.
Experts disagree over whether temnospondyls were ancestral to modern amphibians (frogs, salamanders and caecilians), or whether the whole group died out without leaving any descendants. Different hypotheses have placed modern amphibians as the descendants of temnospondyls, as descendants of another group of early tetrapods called lepospondyls, or even as descendants of both groups (with caecilians evolving from lepospondyls and frogs and salamanders evolving from temnospondyls). There is further disagreement about a temnospondyl origin of lissamphibians related to whether the modern groups arose from only one group (dissorophoids) or from two different groups (dissorophoids and stereospondyls). The majority of studies place a group of temnospondyls called amphibamiforms as the closest relatives of modern amphibians. Similarities in teeth, skulls and hearing structures link the two groups. Whether temnospondyls are considered part of the tetrapod crown or stem thus depends on their inferred relationship to lissamphibians.
Definitions
Branch-based definition
In 2000, Adam Yates and Anne Warren defined the name Temnospondyli as applying to the clade encompassing all organisms that are more closely related to Eryops than to the “microsaur” Pantylus. By this definition, if lissamphibians are temnospondyls and Pantylus is a reptiliomorph, the name Temnospondyli is synonymous with Batrachomorpha (a clade containing all organisms that are more closely related to modern amphibians than to mammals and reptiles).
Node-based definition
Rainer Schoch in 2013 defined the name Temnospondyli as applying to “[t]he least inclusive clade containing Edops craigi and Mastodonsaurus giganteus”.
Description
Many temnospondyls are much larger than living amphibians, and superficially resemble crocodiles, which has led many taxa to be named with the suffix -suchus. The largest taxa, which were predominantly the Mesozoic stereospondyls, had skulls exceeding one meter in length, and the entire animal would have been several meters in length (for reference, the largest living amphibian, Andrias, is about 1.8 meters in body length). Others are smaller and resemble salamanders, in particularly the amphibamiform and micromelerpetid dissorophoids.
Cranium
Skulls are generally parabolic to triangular in shape when viewed from above, and they were particularly flattened in semiaquatic to aquatic taxa, with dorsally facing orbits. The skull is usually covered in pits and ridges to form a honeycomb-like pattern. One of the most recent hypotheses for the function of the dermal ornamentation is that it may have supported blood vessels, which could transfer carbon dioxide to the bones to neutralize acidic build up in the blood (early semiaquatic tetrapods would have had difficulty expelling carbon dioxide from their bodies while on land, and these dermal bones may have been an early solution to the problem). However, there are many other possible hypotheses for the purpose of the ornamentation (e.g., increasing surface area for better adhesion of the skin to the skull), and the function(s) remains largely unresolved due to the absence of this feature in lissamphibians. Some temnospondyls also exhibit raised tubercles or pustules instead of pits and grooves (e.g., the dissorophoid Micropholis, plagiosaurine plagiosaurids), and the import of this disparity is also unclear. Many temnospondyls also have canal-like grooves in their skulls called sensory sulci, the presence of which is used to infer an aquatically inclined lifestyle. The sulci, which usually run around the nostrils and eye sockets, are part of a lateral line system used to detect vibrations in water in modern fish and certain modern amphibians. Many taxa, especially those inferred to have been terrestrial, have an opening at the midline near the tip of the snout called the internarial fenestra / fontanelle; this may have housed a mucous gland used in prey capture. In zatracheids, this opening is greatly enlarged for an unknown purpose.
Homologues of most of the bones of temnospondyls are also seen in other early tetrapods, aside from a few bones in the skull, such as interfrontals, internasals and interparietals, that have developed in some temnospondyl taxa. The intertemporal, a bone common in stem tetrapods, is only found in some late Paleozoic taxa like certain edopoids and dvinosaurs. Most temnospondyls have an indentation at the back of the skull called otic notches. It has typically been inferred that this structure supported a typanum for hearing, although there is substantial variation among temnospondyls in the anatomy of this notch such that it may not have served this function in all temnospondyls, and some clades like plagiosaurids and brachyopids lack notches entirely.
The palate of temnospondyls generally consists of the same bones found in other early tetrapods. Among the most distinguishing features of temnospondyls are the interpterygoid vacuities, two large holes in the back of the palate. Recent studies have suggested that these large openings provided additional attachment sites for musculature and that many temnospondyls were capable of retracting their eyeballs through the vacuities, which is observed in modern frogs and salamanders that also have these large palatal openings; there is no evidence for a buccal pump mechanism for respiration. Temnospondyls often have extensive coverings of teeth on their palates, as well as in their jaws, in contrast to modern amphibians. Some of these teeth are so large that they are referred to as tusks or fangs. Although most temnospondyls have monocuspid teeth, the presence of bicuspid and/or pedicellate teeth in some dissorophoids has been cited as evidence for close relatedness to lissamphibians. In some temnospondyls, such as the dvinosaur Erpetosaurus, the capitosaur Mastodonsaurus and the trematosaur Microposaurus, tusks in the lower jaw pierce the palate and emerge through openings in the top of the skull.
Postcranium
Temnospondyls' vertebrae are divided into several segments. In living tetrapods, the main body of the vertebra is a single piece of bone called the centrum, but in temnospondyls, this region was divided into a pleurocentrum and intercentrum. Two primary types of vertebrae are recognized in temnospondyls: stereospondylous and rhachitomous vertebrae. In rhachitomous vertebrae, the intercentra are large and wedge-shaped, and the pleurocentra are relatively small blocks that fit between them. Both elements support a spine-like neural arch, and well-developed interlocking projections called zygapophyses strengthen the connections between vertebrae. The strong backbone and strong limbs of many rhachitomous temnospondyls allowed them to be partially, and in some cases fully, terrestrial. In stereospondylous vertebrae, the pleurocentra have been greatly reduced or lost entirely, with the intercentra enlarged as the main body of the vertebrae. Early concepts of stereospondyl required the pleurocentra to be entirely absent, but newer concepts only require that the intercentrum has become greatly enlarged. This weaker type of backbone indicates that stereospondylous temnospondyls spent more time in water. Additional types that are less common are the plagiosaurid-type in which there is a single enlarged centrum of uncertain homology; and the tupilakosaurid-type vertebrae (diplospondyly) in which the pleurocentra and intercentra are the same size and form discs; this occurs in tupilakosaurid dvinosaurs but also at least some brachyopids and several other non-temnospondyls. The neural spines tend to be of similar height throughout the presacral region of the trunk, but some temnospondyls exhibit increasing height towards the mid-trunk, followed by a decrease in height to produce a more hump-backed contour. The most extreme is observed in the dissorophid Platyhystrix, which has greatly elongated neural spines that form a large sail on its back. The function of this sail, like that of the contemporaneous sphenacodontids and edaphosaurids, remains enigmatic, but it is thought to have stiffened the vertebral column in association with the relative terrestriality of this clade. Recent histological work has demonstrated that most of this hyperelongation is formed by the osteoderm capping the spine, and thus the sail of Platyhystrix is dissimilar to that of pelycosaurs in which it is entirely formed by the spine. The majority of temnospondyls have presacral counts between 23 and 27, with reduction observed in some amphibamiforms and elongation observed in many dvinosaurs. Caudal length is highly variable, and complete caudal sequences are rare. Based on Eryops, more than 30 caudal positions were possible in some taxa.
The pectoral girdle comprised an unpaired interclavicle, paired clavicles, paired cleithra, and paired scapulae / scapulocoracoids as with most other early tetrapods. These elements differ widely in variation across temnospondyls, with such variation attributed to different lifestyles. The interclavicle and clavicles tend to be more lightly built in terrestrial taxa, with little to no ornamentation. In contrast, these elements are massively ossified in the aquatic stereospondyls and are well ornamented in the same fashion as the skull. The cleithrum and scapulocoracoid is more developed in terrestrial taxa, and the coracoid tends not to ossify in aquatic forms such that there is only a much shorter scapula present. The pelvis comprises the ilium, ischium and pubis, the last of which does not always ossify in aquatic forms. The sutural contacts between elements may also be visible, even when all three ossify. The forelimb comprised the typical radius, ulna, humerus and manus. These bones are typically more developed with greater surface area for muscle attachment in taxa inferred to have been terrestrial. Many dissorophoids have long and slender limbs. Historically it has been thought that all temnospondyls had only four fingers, but this has been shown not to be true in at least a few stereospondyls (Metoposaurus, Paracyclotosaurus), and the paucity of complete manuses casts doubt on the sweeping characterization of a four-fingered manus as the predominant or plesiomorphic condition. At least in Metoposauridae, there are both taxa with four fingers and taxa with five. The hindlimb comprised the typical tibia, fibula, femur and pes. Relative development is as with the forelimb. All temnospondyls with a known pes have five digits.
Unlike modern amphibians, many temnospondyls are covered in small, closely packed scales. The undersides of most temnospondyls are covered in rows of large ventral plates. During early stages of development, they first have only small, rounded scales. Fossils show, as the animals grew, the scales on the undersides of their bodies developed into large, wide ventral plates. The plates overlap each other in a way that allows a wide range of flexibility. Later semiaquatic temnospondyls, such as trematosaurs and capitosaurs, have no evidence of scales. They may have lost scales to make movement easier under water or to allow cutaneous respiration, the absorption of oxygen through the skin.
Several groups of temnospondyls have large bony plates (osteoderms) on their backs. One temnospondyl, Peltobatrachus, has armour-like plating that covers both its back and underside. The rhytidosteid Laidleria also has extensive plating on its back. Most members of the family Dissorophidae also have armor, although it only covers the midline of the back with one or two narrow rows of plates that tightly articulated with the vertebrae, and osteoderms are also known from a few trematopids. Other temnospondyls, such as Eryops, have been found with small, disc-like bony scutes that were in life probably embedded in the skin. All of these temnospondyls were adapted to a terrestrial lifestyle. Armor may have offered protection from predators in the case of Peltobatrachus. The scutes may have provided stability for the spine, as they would have limited flexibility and may have been connected by strong ligaments. A carapace of osteoderms is also seen in plagiosaurids, the only primarily aquatic clade with such extensive ossifications. Plagiosaurids may have inherited their armor from a terrestrial ancestor, as both Peltobatrachus and Laidleria have been considered close relatives of the group. Alternatively, these osteoderms may have served as mineral reservoirs to allow plagiosaurids to respond to a variety of environmental conditions. Contrary to older assumptions, more recent studies have argued that the temnospondyls evolved from a terrestrial ancestor (although with aquatic eggs and larvae), and that it was the forms that later returned to water and an aquatic lifestyle which evolved a spine more rigid and stiffer than the terrestrial species.
Soft tissue
Very little is known of the soft tissue of temnospondyls because the conditions necessary to preserve such material are uncommon. The most extensive records come from fine-grained deposits in the Carboniferous and Permian of Germany; the small-bodied and aquatic dissorophoids and the larger stereospondylomorphs are frequently preserved with outlines of soft tissue around the skeleton. Typically preserved features include the outline of the body, external gills, and parts of the eye or stomach. An amphibamiform specimen from the Mazon Creek locality was described as having toepad-like features. The holotype specimen of Arenaerpeton supinatus from the Triassic of New South Wales, Australia, displays extensive soft tissue, hinting at the girth of the animal in life. Trace fossils attributed to temnospondyls are fairly common, especially from the Carboniferous through the Triassic. Common ichnogenera include Batrachichnus and Limnopus.
History of study
Temnospondyli was named by the German paleontologist Karl Alfred von Zittel in his second edition of Handbuch der Palaeontologie, published in 1888. However, temnospondyl remains have been known since the early part of the 19th century.
Early finds: Mastodonsaurus and "labyrinthodonts" (early to mid-19th century)
The earliest described temnospondyl was Mastodonsaurus, named by Georg Friedrich Jaeger in 1828 from a single tooth that he considered to belong to a reptile. Mastodonsaurus means "breast tooth lizard" after the nipple-like shape of the tip of the tooth.
The naming of these first specimens was disputed. Leopold Fitzinger named the animal Batrachosaurus in 1837. In 1841, the English paleontologist Richard Owen referred to the genus as Labyrinthodon to describe its highly folded or labyrinthine teeth. Owen thought that the name Mastodonsaurus "ought not to be retained, because it recalls unavoidably the idea of the mammalian genus Mastodon, or else a mammilloid form of the tooth... and because the second element of the word, saurus, indicates a false affinity, the remains belonging, not to the Saurian, but to the Batrachian order of Reptiles." Owen recognized that the animal was not a "saurian" reptile, yet he also referred Jaeger's Phytosaurus to the genus. Although the two genera have similarly sized conical teeth, Phytosaurus was later found to be a crocodile-like reptile. Additional material, including skulls, firmly placed Labyrinthodon as an amphibian. Jaeger also named Salamandroides giganteus in 1828, basing it on partial occiput, or back portion of the skull. In 1833, he described a complete skull of S. giganteus that had the same teeth as his Mastodonsaurus, making it the first-known complete skull of a temnospondyl. Because Mastodonsaurus was named first, it has precedence over the other names as a senior subjective synonym.
Mastodonsaurus and other similar animals were referred to as labyrinthodonts, named like Labyrinthodon for teeth that were highly folded in cross section. Owen's "Labyrinthodon Jaegeri" was later found at Guy's Cliffe, England by paleontologist William Buckland. Other specimens were found in the red sandstone of Warwickshire. As more fossils were uncovered in England, Owen depicted these labyrinthodonts as the "highest" form of batrachian and compared them to crocodiles, which he considered the highest form of reptiles. He also noted the large labyrinthodonts of the Keuper (a unit of rocks that dates to the Late Triassic) were younger than more advanced reptiles in the Magnesian and Zechstein, which are Late Permian in age. Owen used these fossils to counter the notion that reptiles evolved from a sequential progression from early amphibians (what he called "metamorphosed fishes").
In addition to Mastodonsaurus, some of the earliest-named genera included Metopias and Rhombopholis in 1842, Zygosaurus in 1848, Trematosaurus in 1849, Baphetes and Dendrerpeton in 1853, Capitosaurus in 1858, and Dasyceps in 1859. Baphetes is now placed as an early tetrapod outside Temnospondyli, and Rhombopholis is now considered a prolacertiform reptile.
Labyrinthodonts as amphibians (late 19th century)
Later in the 19th century, temnospondyls were classified as various members of Stegocephalia, a name coined by the American paleontologist Edward Drinker Cope in 1868. Cope placed stegocephalians in the class Batrachia, the name then used for Amphibia. Stegocephalia means "roof-headed" in Greek, a reference to the wide, flat heads of temnospondyls and other early tetrapods. During this time, paleontologists considered temnospondyls to be amphibians because they possessed three main features: gill arches in juvenile skeletons, indicating they were amphibious for at least the first part of their lives; ribs that do not connect at the underside of the rib cage; and deep pits in the skull that were interpreted as space for mucous glands.
Several suborders of stegocephalians were recognized in the late 19th and early 20th centuries. Animals now regarded as temnospondyls were primarily labyrinthodonts, but some were classified in the Branchiosauria. Branchiosaurs were small-bodied and had simple conical teeth, while labyrinthodonts were larger and had complex, folded dentin and enamel in their teeth. Branchiosauria included only a few forms, such as Branchiosaurus from Europe and Amphibamus from North America, that had poorly developed bones, external gills, and no ribs. Some skeletons of Amphibamus were later found with long ribs, prompting its reassignment to Microsauria (although more detailed studies found it to be a temnospondyl). Soft tissue, such as scales and external gills, were found in many well-preserved branchiosaur fossils from Germany. In the early 20th century, branchiosaurs would be recognized as larval forms of temnospondyls lacking many of the typical features that define the group, and is no longer recognized as a distinct group.
Other animals that would later be classified as temnospondyls were placed in a group called Ganocephala, which was characterized by plate-like skull bones, small limbs, fish-like scales and branchial arches. Unlike labyrinthodonts, they did not have parietal foramina, small holes in their skulls behind their eye sockets. Archegosaurus, Dendrerpeton, Eryops and Trimerorhachis were placed in this group and were considered to be the most primitive members of Reptilia. Their rhachitomous vertebrae, notochord and lack of occipital condyles (which attached the head to the neck) were features that were also shared with fishes. Thus, they were considered a link between early fishes and more advanced forms such as stegocephalians.
Another group was called Microsauria by Cope in 1868. He classified Microsauria as a subgroup of Labyrinthodontia, placing many small, amphibian-like animals within it. Among them was Dendrerpeton, once placed in Ganocephala. Dendrerpeton was later placed as a labyrinthodont with other temnospondyls, but confusion existed for many years over the classification of small amphibians.
By the end of the 19th century, most of what are today regarded as temnospondyls were placed in the suborder Labyrinthodonta. The American paleontologist Ermine Cowles Case called it Labyrinthodonta vera or "true labyrinthodonts". The names Stegocephalia and Labyrinthodontia were used interchangeably to refer to the order in which it belonged. The labyrinthodontian suborders Microsauria and Branchiosauria, both of which contain temnospondyls, were distinct from Labyrinthodonta. Within Labyrinthodonta were the groups Rhachitomi, Labyrinthodonti and Embolerimi. Members of Rhachitomi, such as Archegosaurus and Eryops, had rhachitomous vertebrae with enlarged intercentra that displaced the pleurocentra. Labyrinthodonti, such as Mastodonsaurus, Trematosaurus and Micropholis, had lost their pleurocentra, and the intercentra made up the entire body of the vertebrae. Embolerimi had intercentra and pleurocentra that were of equal size. Embolomeres are now identified as a separate group of reptiliomorphs or stem-group tetrapods, with no particular affinities to temnospondyls.
Vertebra-based classifications and the origin of the name "Temnospondyli" (1888 – 20th century)
In 1888, von Zittel divided stegocephalians among three taxa: Lepospondyli, Temnospondyli and Stereospondyli. He placed microsaurs in Lepospondyli, a group which he characterized as having simple, spool-shaped vertebral centra. Temnospondyli included forms with the centra divided into pleurocentra and intercentra. All members of Stereospondyli had amphicoelous centra composed only of the intercentra. Cope objected to von Zittel's classification, considering the vertebrae of lepospondyls and stereospondyls indistinguishable because each had a simple spool shape. He continued to use Ganocephala and Labyrinthodonta (which he alternatively referred to as Rhachitomi) to distinguish animals based on the absence or presence of occipital condyles.
Temnospondyli became a commonly used name at the turn of the 20th century. Paleontologists included both embolomeres and rhachitomes in the group. Cope's Ganocephala and Labyrinthodonta fell out of use. In 1919, British paleontologist D. M. S. Watson proposed that the evolutionary history of these large amphibians could be seen through changes in their vertebrae. Embolomerous forms in the Carboniferous graded into rhachitomous forms in the Permian, and finally into stereospondyls in the Triassic. More importantly, Watson began using the term Labyrinthodontia to refer to these groups. The name Temnospondyli was rarely used in the decades that followed. Swedish paleontologist Gunnar Säve-Söderbergh removed embolomeres from the group, narrowing its scope to rhachitomes and stereospondyls. His classification of labyrinthodonts was based heavily on characteristics of the skull rather than the vertebrae.
The American paleontologist Alfred Romer brought the name Temnospondyli back into use in the later 20th century. Säve-Söderbergh used the name Labyrinthodontia in a strict sense (sensu stricto) to refer to Rhachitomi and Stereospondyli, excluding Embolomeri. Romer agreed with this classification, but used the name Temnospondyli to avoid confusion with Labyrinthodontia in its wider sense (sensu lato). Unlike modern temnospondyl classification, however, Romer included the primitive Ichthyostegalia in the group.
20th century – present
More recent study of temnospondyls has largely focused on their paleobiology and resolving their internal relationships. With a few exceptions, the monophyly of Temnospondyli is not questioned, and the disuse of terms like Labyrinthodontia and Stegocephalia continues. Temnospondyls continue to be heavily involved in the debate over lissamphibian origins.
As with evolutionary biology in general, computer-assisted phylogenetic methods have greatly facilitated phylogenetic inference of the relationships of both Temnospondyli at large and specific sub-groups. Other quantitative analyses have addressed morphometrics, biomechanics,
Temnospondyls were also documented from an increasingly broad geographic and stratigraphic range in the 20th and 21st centuries, including the first occurrences from historically undersampled regions such as Antarctica, Lesotho, Japan, Namibia, New Zealand, Niger, and Türkiye.
Evolutionary history
Carboniferous and Early Permian (Cisuralian)
Temnospondyls first appeared in the Middle Mississippean (Viséan) around 330 million years ago (Mya) where the earliest appearances are Balanerpeton from Scotland and an indeterminate temnospondyl from Germany. During the Carboniferous, all of the non-stereospondylomorph clades appeared, including dendrerpetids, edopoids, eryopoids, the various dissorophoid subclades, dvinosaurs and zatracheids. Stereospondylomorphs and stereospondyls first appeared in the early Permian, although the former may have appeared earlier and merely be undocumented at present. The vast majority of the Carboniferous records come from the midwestern United States, such as the Linton, Five Points and Mazon Creek lagerstätte, and the south-central United States where classic redbed formations are found; and from western Europe, particularly the Saar-Nahe Basin in Germany and Nýřany in the Czech Republic. The early Permian record of temnospondyls is also concentrated in these regions. Most of the clades from the Late Carboniferous continued to be successful, with a particularly high diversity of dissorophoids.
Middle Permian (Guadalupian)
Middle Permian records of temnospondyls are relatively sparse, and some of these are debated as a result of the uncertain age and correlation of different deposits in North America (Chickasha, Flowerpot Formations), Niger (Moradi Formation), Brazil (Rio do Rasto Formation), and Russia (Mezen complex) and the subsequent controversy over Olson's Gap, but there is a general consensus that at least some of these records are Guadalupian in age. Records of rhinesuchids from the Eodicynodon andTapinocephalus Assemblage Zones of South Africa are less controversial. Additional records are known from Brazil, China, Turkey, and the Isheevo complex of Russia. A mixture of taxa are represented, including stereospondylomorphs (Konzhukovia) and rhinesuchid stereospondyls, as well as some of the latest occurrences of dissorophoids (Anakamacops, Kamacops).
Late Permian (Lopingian)
During the late Permian, increasing aridity and the diversification of reptiles contributed to a decline in terrestrial temnospondyls, but semiaquatic and fully aquatic stereospondylomorph temnospondyls continued to flourish, including the large Melosaurus of Eastern Europe. Other temnospondyls, such as archegosaurids, developed long snouts and a close similarity to crocodiles, although they lacked the armor characteristic of the latter group. These temnospondyls included the largest-known batrachomorph, the over 5.5-meter-long Prionosuchus of Brazil. The stereospondyl record is almost exclusively confined to rhinesuchids.
Triassic
As temnospondyls continued to flourish and diversify in the Late Permian (260.4–251.0 Mya), a major group called Stereospondyli became more dependent on life in the water. The vertebrae became weak, the limbs small, and the skull large and flat, with the eyes facing upwards. During the Triassic period, these animals dominated the freshwater ecosystems, evolving in a range of both small and large forms. During the Early Triassic (251.0–245.0 Mya) one group of successful long-snouted fish-eaters, the trematosauroids, even adapted to a life in the sea, the only known batrachomorphs to do so with the exception of the modern crab-eating frog. Another group, the capitosauroids, included medium-sized and large animals in length, with large and flat skulls that could be over a meter long in the largest forms such as Mastodonsaurus. These animals spent most or all their lives in water as aquatic predators, catching their prey by a sudden opening of the upper jaw and sucking in fish or other small animals.
In the Carnian stage of the Late Triassic (237.0–227.0 Mya), capitosauroids were joined by the superficially very similar Metoposauridae. Metoposaurids are distinguished from capitosauroids by the positioning of their eye sockets near the front of their skulls. Another group of stereospondyls, the plagiosaurs, had wide heads and gills, and adapted to life at the bottom of lakes and rivers. By this time, temnospondyls had become a common and widespread component of semiaquatic ecosystems. Some temnospondyls, such as Cryobatrachus and Kryostega, even inhabited Antarctica, which was covered in temperate forests at the time.
Triassic temnospondyls were often the dominant semiaquatic animals in their environments. Large assemblages of Late Triassic metoposaurids with hundreds of individuals preserved together have been found in the southwestern United States, Morocco, India, and western Europe. They have often been interpreted as mass death events caused by droughts in floodplain environments. Recent studies show these dense assemblages were instead probably the result of currents transporting and accumulating dead individuals in certain areas.
Jurassic and Cretaceous
Temnospondyls reached a peak diversity during the Early Triassic, and progressively declined throughout the subsequent Middle and Late Triassic, with only members of the Brachyopoidea and Trematosauroidea surviving into the Jurassic and the Cretaceous. Among brachyopoids, the brachyopids Gobiops and Sinobrachyops are known from Middle and late Jurassic deposits across Asia and the chigutisaurid Siderops is known from the Early Jurassic of Australia. The most recent known temnospondyl was the giant chigutisaurid Koolasuchus, known from the Early Cretaceous (Aptian) of Australia. It is thought to have survived in rift valleys that were too cold in the winter for crocodylomorphs that normally would have competed with them. Koolasuchus was one of the largest of the brachyopoids, with an estimated weight of .
Classification
Originally, temnospondyls were classified according to the structure of their vertebrae. Early forms, with complex vertebrae consisting of a number of separate elements, were placed in the suborder Rachitomi, and large Triassic aquatic forms with simpler vertebrae were placed in the suborder Stereospondyli. With the recent growth of phylogenetics, this classification is no longer viable. The basic rhachitomous condition is found in many primitive tetrapods, and is not unique to one group of temnospondyls. Moreover, the distinction between rhachitomous and stereospondylous vertebrae is not entirely clear. Some temnospondyls have rhachitomous, semirhachitomous and sterospondylous vertebrae at different points in the same vertebral column. Other taxa have intermediate morphologies that do not fit into any category. Rachitomi is no longer recognized as an exclusive group (i.e. it includes Stereospondyli rather than being a counterpart to it), but Stereospondyli is still considered valid. Below is a simplified taxonomy of temnospondyls showing currently recognized groups:
Class Amphibia
Order Temnospondyli
Superfamily Edopoidea
Family Cochleosauridae
Family Edopidae
Clade Eutemnospondyli
Family Dendrerpetidae
Clade Rhachitomi
Suborder Dvinosauria
Family Trimerorhachidae
Superfamily Dvinosauroidea
Family Dvinosauridae
Family Eobrachyopidae
Family Tupilakosauridae
Superfamily Dissorophoidea
Family Micromelerpetontidae
Clade Xerodromes
Clade Amphibamiformes
Family Amphibamidae
Family Branchiosauridae
Subfamily Branchiosaurinae
Family Micropholidae
Subclass Lissamphibia? (placement is uncertain)
Clade Olsoniformes
Family Dissorophidae
Subfamily Dissorophinae
Subfamily Eucacopinae
Family Trematopidae
Family Zatracheidae
Clade Eryopiformes
Suborder Euskelia
Superfamily Dissorophoidea?
Superfamily Eryopoidea
Family Eryopidae
Family Zatracheidae?
Clade Limnarchia
Suborder Dvinosauria?
Clade Stereospondylomorpha
Superfamily Archegosauroidea
Family Archegosauridae
Subfamily Melosaurinae
Subfamily Platyoposaurinae
Family Intasuchidae?
Family Sclerocephalidae
Suborder Stereospondyli
Family Peltobatrachidae
Family Lapillopsidae?
Family Rhinesuchidae
Clade Superstes
Family Lydekkerinidae
Clade Neostereospondyli
Clade Capitosauria
Family Sclerothoracidae
Superfamily Mastodonsauroidea
Family Heylerosauridae
Family Mastodonsauridae
Family Stenotosauridae
Clade Trematosauria
Superfamily Trematosauroidea
Family Thoosuchidae
Family Trematosauridae
Subfamily Tertreminae
Subfamily Lonchorhynchinae
Subfamily Trematosaurinae
Superfamily Metoposauroidea
Family Latiscopidae
Family Metoposauridae
Superfamily Plagiosauroidea
Family Laidleriidae?
Family Plagiosauridae
Superfamily Rhytidosteoidea
Family Indobrachyopidae?
Family Rhytidosteidae
Subfamily Derwentiinae
Clade Brachyopomorpha
Superfamily Brachyopoidea
Family Brachyopidae
Family Chigutisauridae
Phylogeny
In one of the earliest phylogenetic analyses of the group, Gardiner (1983) recognized five characteristics that made Temnospondyli a clade: a bone at the back of the skull, the parasphenoid, is connected to another bone on the underside of the skull, the pterygoid; large openings called interpterygoid vacuities are present between the pterygoids; the stapes (a bone involved in hearing) is connected to the parasphenoid and projects upward; the cleithrum, a bone in the pectoral girdle, is thin; and part of the vertebra called the interdorsal attaches to the neural arch. Additional features were given by Godfrey et al. (1987), including the contact between the postparietal and exoccipital at the back of the skull, small projections (uncinate processes) on the ribs, and a pelvic girdle with each side having a single iliac blade. These shared derived characteristics are called synapomorphies.
Temnospondyls are placed as basal tetrapods in phylogenetic analyses, with their exact positioning varying between studies. Depending on the classification of modern amphibians, they are either included in the crown group Tetrapoda or the stem of Tetrapoda. Crown-group tetrapods are descendants of the most recent common ancestor of all living tetrapods and stem tetrapods are forms that are outside the crown group. Modern amphibians have recently been suggested as descendants of temnospondyls, which would place them within crown Tetrapoda. Below is a cladogram from Ruta et al. (2003) placing Temnospondyli within crown Tetrapoda:
Other studies place modern amphibians as the descendants of lepospondyls and place temnospondyls in a more basal position within the stem of Tetrapoda. Below is a cladogram from Laurin and Reisz (1999) placing Temnospondyli outside crown Tetrapoda:
Most phylogenetic analyses of temnospondyl interrelationships focus on individual families. One of the first broad-scale studies of temnospondyl phylogeny was conducted by paleontologist Andrew Milner in 1990. A 2007 study made a "supertree" of all temnospondyl families, combining the family-level trees of previous studies. The following cladogram is modified from Ruta et al. (2007):
1 Temnospondyli, 2 Edopoidea, 3 Dvinosauria, 4 Euskelia, 5 Eryopoidea, 6 Dissorophoidea, 7 Limnarchia, 8 Archegosauroidea,
9 Stereospondyli, 10 Rhytidostea, 11 Brachyopoidea, 12 Capitosauria, 13 Trematosauria, 14 Metoposauroidea
The most basal group of temnospondyls is the superfamily Edopoidea. Edopoids have several primitive or plesiomorphic features, including a single occipital condyle and a bone called the intertemporal that is absent in other temnospondyls. Edopoids include the Late Carboniferous genus Edops and the family Cochleosauridae. Dendrerpetontidae has also been included in Edopoidea, and is the oldest-known temnospondyl family. Balanerpeton woodi is the oldest species, having been present over 330 million years ago during the Viséan stage of the Early Carboniferous. Recent analyses place Dendrerpetontidae outside Edopoidea in a more derived position. Other primitive temnospondyls include Capetus and Iberospondylus. Saharastega and Nigerpeton, both described in 2005 from Niger, are also primitive yet come from the Late Permian. They are almost 40 million years younger than other basal temnospondyls, implying a long ghost lineage of species that are not yet known in the fossil record.
In 2000, paleontologists Adam Yates and Anne Warren produced a revised phylogeny of more derived temnospondyls, naming several new clades. Two major clades were Euskelia and Limnarchia. Euskelia includes the temnospondyls that were once called rhachitomes and includes two subfamilies, the Dissorophoidea and the Eryopoidea. Dissorophoids include small, mostly terrestrial temnospondyls that may be the ancestors of modern amphibians. Eryopoids include larger temnospondyls like Eryops. The second major clade, Limnarchia, includes most Mesozoic temnospondyls, as well as some Permian groups. Within Limnarchia are the superfamily Archegosauroidea and the most derived temnospondyls, the stereospondyls.
Yates and Warren also named Dvinosauria, a clade of small aquatic temnospondyls from the Carboniferous, Permian and Triassic. They placed Dvinosauria within Limnarchia, but more recent studies disagree on their position. For example, a 2007 study places them even more basal than euskelians, while a 2008 study keeps them as basal limnarchians.
Within Stereospondyli, Yates and Warren erected two major clades: Capitosauria and Trematosauria. Capitosaurs include large semiaquatic temnospondyls like Mastodonsaurus with flat heads and eyes near the back of the skull. Trematosaurs include a diversity of temnospondyls, including large marine trematosaurids, aquatic plagiosaurs, brachyopoids that survived into the Cretaceous, and metoposauroids with eyes near the front of their heads. In 2000, paleontologists Rainer Schoch and Andrew Milner named a third major clade of stereospondyls, the Rhytidostea. This group included more primitive stereospondyls that could not be placed in either Capitosauria or Trematosauria, and included groups like Lydekkerinidae, Rhytidosteidae and Brachyopoidea. While Capitosauria and Trematosauria are still widely used, Rhytidostea is not often supported as a true clade in recent analyses. Rhytidosteids and brachyopoids are now grouped with trematosaurians, but lydekkerinids are still considered to be a primitive family of stereospondyls.
A new phylogeny of temnospondyls was offered by paleontologist Rainer Schoch in 2013. It supported many of the clades that were found by Yates and Warren, but it did not find support for their division of derived stereospondyls into Euskelia and Limnarchia. Eryopids were found to be more closely related to stereospondyls than to dissorophoids, which were grouped with dvinosaurs. The clade including Eryopidae and Stereospondylomorpha was named Eryopiformes. In addition, Schoch named the clade containing all temnospondyls except edopoids Eutemnospondyli and reinstated the name Rhachitomi for the clade containing all temnospondyls except edopoids and dendrerpetontids. Below is the cladogram from Schoch's analysis:
Relationship to modern amphibians
Modern amphibians (frogs, salamanders and caecilians) are classified in Lissamphibia. Lissamphibians appear to have arisen in the Permian. Molecular clock estimates place the first lissamphibian in the Late Carboniferous, but the first member of Batrachia (frogs and salamanders, but not caecilians) is estimated to have appeared in the Middle Permian using the same technique.
Using fossil evidence, there are three main theories for the origin of modern amphibians. One is that they evolved from dissorophoid temnospondyls. Another is that they evolved from lepospondyls, most likely the lysorophians. A third hypothesis is that caecilians descended from lepospondyls and frogs and salamanders evolved from dissorophoids.
Recently, the theory that temnospondyls were the ancestors of all lissamphibians has gained wide support. The skull morphology of some small temnospondyls has been compared to those of modern frogs and salamanders, but the presence of bicuspid, pedicellate teeth in small, paedomorphic or immature temnospondyls has been cited as the most convincing argument in favor of the temnospondyl origin of lissamphibians. Seen in lissamphibians and many dissorophoid temnospondyls, pedicellate teeth have calcified tips and bases. During the development of most tetrapods, teeth begin to calcify at their tips. Calcification normally proceeds downward to the base of the tooth, but calcification from the tip stops abruptly in pedicellate teeth. Calcification resumes at the base, leaving an area in the center of the tooth uncalcified. This pattern is apparent in both living amphibians and certain dissorophoid fossils.
The dissorophoid family Amphibamidae is thought to be most closely related to Lissamphibia. In 2008, an amphibamid called Gerobatrachus hottoni was named from Texas and was nicknamed the "frogamander" for its frog-like head and salamander-like body. It was thought to be the most closely related temnospondyl to lissamphibians and was placed as the sister taxon of the group in a phylogenetic analysis. Another species of amphibamid called Doleserpeton annectens is now thought to be even more closely related to lissamphibians. Unlike Gerobatrachus, Doleserpeton was known since 1969, and the presence of pedicellate teeth in its jaws has led some paleontologists to conclude soon after its naming that it was a relative of modern amphibians. It was first described as a "protolissamphibian", and the specific name annectens means "connecting" in reference to its inferred transitional position between temnospondyls and lissamphibians. The structure of its tympanum, a disk-like membrane that functions like an ear drum, is similar to that of frogs and has also been used as evidence for a close relationship. Other features including the shape of the palate and the back of the skull, the short ribs, and the smooth skull surface also point to it being a closer relative of lissamphibians than is Gerobatrachus. Below is a cladogram modified from Sigurdsen and Bolt (2010) showing the relationships of Gerobatrachus, Doleserpeton and Lissamphibia:
Chinlestegophis, a putative Triassic stereospondyl considered to be related to metoposauroids such as Rileymillerus, has been noted to share many features with caecilians, a living group of legless burrowing amphibians. If Chinlestegophis is indeed both an advanced stereospondyl and a relative of caecilians, this means that although all lissamphibians are descended from temnospondyls, the different groups would have descended from different branches of the temnospondyl family tree. Anurans and urodelans would therefore be surviving dissorophoids while apodans (caecilians) are surviving stereospondyls.
Paleobiology
Metabolism and gas exchange
A study on the fully aquatic Archegosaurus shows that its heat balance, gas exchange, osmoregulation, and digestion were more similar to those of fish than those of modern aquatic amphibians like salamanders.
Feeding
Although the earliest temnospondyls were primarily semiaquatic, they had the ability to feed on land. Later, eryopoids and dissorophoids, some well adapted to terrestrial life, also fed on land. Some eryopoids became better adapted toward life in water and shifted their diets toward aquatic organisms. The first primarily aquatic feeders were archegosaurs in the Permian. Trematosaurs and capitosaurs became independently aquatic and also returned to this type of feeding.
Most aquatic stereospondyls have flattened heads. When feeding, they probably opened their mouths by lifting their skulls instead of lowering their lower jaws. The jaw mechanics of the plagiosaurid Gerrothorax is well known, and is one of the most highly adapted. Gerrothorax is thought to have lifted its skull to around 50 degrees above horizontal through the flexing of the atlanto-occipital joint between the occipital condyles of the skull and the atlas vertebra of the neck. As the skull is raised, the quadrate bone pushes forward and causes the lower jaw to protrude outward. Other stereospondyls probably also lifted their skulls, but they are not as well adapted for such movement. D.M.S. Watson was the first to suggest skull lifting as a means of feeding in temnospondyls. He envisioned that Mastodonsaurus, a much larger temnospondyl than Gerrothorax, was able to make the same movement. Paleontologist A.L. Panchen also supported the idea in 1959, suggesting that Batrachosuchus also fed in this way. At the time it was thought that these temnospondyls lifted their heads with strong jaw muscles, but it is now thought that they used larger muscles in the neck that were attached to the large pectoral girdle. Plagiosuchus, a close relative of Gerrothorax, also has a hyobranchial skeleton that muscles may have attached to. Plagiosuchus has very small teeth and a large area for muscle attachment behind the skull, suggesting that it could suction feed by rapidly opening its mouth.
Unlike semiaquatic temnospondyls, terrestrial temnospondyls have skulls that are adapted for biting land-living prey. The sutures between the bones of the skull in the dissorophoid Phonerpeton are able to withstand a high degree of compression. Compressive forces would have been experienced when biting down on prey. Earlier aquatic tetrapods and tetrapod ancestors differ from temnospondyls like Phonerpeton in that their skulls were also built to withstand tension. This tension would have been experienced during suction feeding underwater. Temnospondyls like Phonerpeton were among the first tetrapods that were almost exclusively terrestrial and fed by biting.
Reproduction
Temnospondyls, like most modern amphibians, reproduced in aquatic environments. Most temnospondyls probably reproduced through external fertilization. Like most living frogs, female temnospondyls would have laid masses of eggs in water while males released sperm to fertilize them. Several fossils were described from the Early Permian of Texas in 1998 that may be egg masses of dissorophoid temnospondyls. They were the first-known fossils of amphibian eggs. The fossils consist of small disks with thin membranes that are probably vitelline membranes and halo-like areas surrounding them that are most likely mucous coatings. They are attached to plant fossils, suggesting that these temnospondyls laid eggs on aquatic plants much like modern frogs. The mucous membranes show that the eggs were laid by amphibians, not fish (their eggs lack mucous), but the type of amphibian that laid them cannot be known because no body fossils are preserved with the eggs. The eggs are thought to be from dissorophoids because they are likely to be close relatives of modern amphibians, and probably had similar reproductive strategies. They are also the most common amphibians from the deposit in which the eggs were found.
One temnospondyl, the dvinosaur Trimerorhachis, may have brooded young in an area between the gills called the pharyngeal pouch. Small bones belonging to younger Trimerorhachis individuals have been found in these pouches. The living Darwin's Frog is also a mouth brooder and would be the closest modern analogue to Trimerorhachis if it cared for its young in this way. An alternative possibility is that Trimerorhachis was cannibalistic, eating its young like many amphibians do today. If this was the case, the bones of these smaller individuals were originally located in the throat and were pushed into the pharyngeal pouch as the animal fossilized.
Body impressions of Early Carboniferous temnospondyls from Pennsylvania suggest that some terrestrial temnospondyls mated on land like some modern amphibians. They reproduced through internal fertilization rather than mating in water. The presence of three individuals in one block of sandstone shows that the temnospondyls were gregarious. The head of one individual rests under the tail of another in what may be a courtship display. Internal fertilization and similar courtship behavior are seen in modern salamanders.
Growth
While most types of temnospondyls are distinguished on the basis of features in mature specimens, several are known from juvenile and larval specimens. Metamorphosis is seen in dissorophoids, eryopids and zatrachydids, with aquatic larvae developing into adults capable of living on land. Several types of dissorophoids, such as branchiosaurids, do not fully metamorphose, but retain features of juveniles such as external gills and small body size in what is known as neoteny. Dvinosaurians and the plagiosaurid Gerrothorax also retained gills, although recent studies found that (at least as adults) their gills were internal like those of fish, rather than external like those of salamanders.
Temnospondyl larvae are often distinguished by poorly developed bones and the presence of a hyobranchial apparatus, a series of bones that gills would attach to in life. However, some fully mature temnospondyls also possess hyobranchial bones but did not have external gills. A dense covering of scales is also seen in larvae and adults. Major body changes occur in metamorphosis, including the reshaping and strengthening of skull bones, the thickening of postcranial bones, and an increase in body size.
Temnospondyls like Sclerocephalus are known from both large adult specimens and small larvae, showing an extreme change in body shape. In these species, the shape and proportions of skull bones change in the early stages of development. The ornamentation on the surface of the skull roof also develops at this time. Small, regularly spaced pits are the first to form, followed by larger ridges. As development continues, the external gills disappear. Small teeth that once covered the palate are lost. The postcranial skeleton does not develop at the same rate as the skull, with ossification (the replacement of cartilage by bone) happening more slowly. Vertebrae and limb bones are poorly developed, ribs and fingers are absent in the early stages, and the scapulocoracoid and ischium are entirely absent through most of development. Once maturity is reached, most bones have fully formed and growth rate slows. The bones of some temnospondyls like Dutuitosaurus show growth marks, possibly an indication that growth rate varied with the change in seasons. Fossils of temnospondyls like Metoposaurus and Cheliderpeton show that individuals grew larger past maturity. The oldest individuals usually have more pitting on their skulls with deeper sulci.
One group of temnospondyls, the Branchiosauridae, is also known from larval specimens. Branchiosaurids like Branchiosaurus and Apateon are represented by many fossils preserving skin and external gills. An entire growth series is exhibited in the wide range of sizes among specimens, but the lack of terrestrially adapted adult forms suggests that these temnospondyls were neotenic. Unlike other temnospondyls, their postcranial skeletons developed quickly but were still partly cartilaginous when fully mature. Adults likely had an aquatic lifestyle similar to juveniles. Recently, large specimens of Apateon gracilis were described with adaptations toward a terrestrial lifestyle, indicating that not all branchiosaurs were neotenic.
Studies of temnospondyl development have reached differing conclusions regarding what forms of gills were present in temnospondyls which possessed the organs. Although some species possessed external gills which were preserved as soft tissue, for many groups the type of gill can only be inferred from the structure of the bones which would have supported them. Scientists have disagreed on what these bones imply. Scientists who compare temnospondyls to fish find that the bones correlate with internal gills, while those who compare them closely to salamanders consider the bones to correlate with external gills. This conundrum, known as Bystrow's paradox, has made it difficult to assess the configuration of gills in aquatic temnospondyls.
Bystrow's paradox was resolved by a 2010 study. This study found that grooved ceratobrachnial structures (components of the branchial arches) are correlated with internal gills. Ancient tetrapods which preserved grooved ceratobranchials, such as the dvinosaur Dvinosaurus, probably only had internal gills as adults. Nevertheless, external gills are known to have been conclusively present in at least some temnospondyls. However, these situations only occur in larval specimens or members of specialized groups such as the branchiosaurids. One living species of lungfish (Lepidosiren) has external gills as larvae which are reconfigured into internal gills as adults. Despite adult Dvinosaurus specimens having skeletal features correlated with internal gills, another dvinosaur, Isodectes, includes larval fossils preserving external gills as soft tissue traces. Thus, the gill development of dvinosaurs (and presumably other temnospondyls) mirrored that of Lepidosiren. Despite this feature likely being an example of convergent evolution (as other lungfish exclusively possessed internal gills), it still remains a useful gauge for how temnospondyl gills developed. The study concluded that fully aquatic gilled temnospondyls (including but not limited to dvinosaurs) possessed internal gills as adults and external gills as larvae.
While most temnospondyls are aquatic in early stages of life, most metoposaurids appear to have been terrestrial in their juvenile stage. Like other Mesozoic temnospondyls, adult metoposaurids were adapted to a semiaquatic lifestyle. Their bones are not highly developed for movement on land. The cross-sectional thickness of limb bones in adult metoposaurids shows that they could not withstand the stress of terrestrial locomotion. Juvenile individuals have bones that are thick enough to withstand this stress, and could probably move about on land. To maintain a terrestrial lifestyle, a temnospondyl's limb bones would have to thicken with positive allometry, meaning that they would grow at a greater rate than the rest of the body. This is not the case in metoposaurids, meaning that as their bodies grew larger they became less adapted toward a terrestrial lifestyle.
Hearing
Temnospondyls and other early tetrapods have rounded otic notches in the back of the skull that project into the cheek region. In life, the otic notch would have been covered by a membrane called the tympanum, which is seen as a disk-like area in living frogs. The tympanum is involved in hearing, and is similar to the ear drum of more advanced tetrapods. It was traditionally thought that the tympanum developed very early in tetrapod evolution as a hearing organ and progressed to form the eardrum of amniotes. Thus, temnospondyls possessed a hearing system supposedly ancestral to that of living amphibians and reptiles.
Frogs and all other living tetrapods have a rod-like bone called the stapes that aids in hearing by transferring vibrations from the ear drum—or homologous tympanum—to the inner ear. Temnospondyls also have a stapes, which projects into the otic cavity. The stapes likely evolved from the hyomandibula of lobe-finned fishes. The positioning of the stapes and the shape of the otic region suggests that the tympani of temnospondyls and frogs are homologous, but the tympani of these amphibians are no longer considered homologous with the hearing systems of reptiles, birds and mammals. Therefore, ear structures in temnospondyls were not ancestral to those of all other tetrapods.
The ability of the tympanum and stapes to effectively transmit vibrations is called impedance matching. Early tetrapods like temnospondyls have thick stapes with poor impedance matching, so it is now thought that they were not used for hearing. Instead, these thick stapes may have functioned to support the tissue that covers the otic notch. Early temnospondyls, like Dendrerpeton, could not hear airborne sound but would have been able to detect vibration in the ground. Later temnospondyls like Doleserpeton had otic regions adapted to hearing. Doleserpeton has a structure in the inner ear called the perilymphatic duct, which is also seen in frogs and is associated with hearing. Its stapes is also a better transmitter of sound. The hearing system of Doleserpeton and related temnospondyls was able to detect airborne sound and may have been ancestral to that of living amphibians.
| Biology and health sciences | Prehistoric amphibians | Animals |
24580536 | https://en.wikipedia.org/wiki/Cobalt | Cobalt | Cobalt is a chemical element; it has symbol Co and atomic number 27. As with nickel, cobalt is found in the Earth's crust only in a chemically combined form, save for small deposits found in alloys of natural meteoric iron. The free element, produced by reductive smelting, is a hard, lustrous, somewhat brittle, gray metal.
Cobalt-based blue pigments (cobalt blue) have been used since antiquity for jewelry and paints, and to impart a distinctive blue tint to glass. The color was long thought to be due to the metal bismuth. Miners had long used the name kobold ore (German for goblin ore) for some of the blue pigment-producing minerals. They were so named because they were poor in known metals and gave off poisonous arsenic-containing fumes when smelted. In 1735, such ores were found to be reducible to a new metal (the first discovered since ancient times), which was ultimately named for the kobold.
Today, some cobalt is produced specifically from one of a number of metallic-lustered ores, such as cobaltite (CoAsS). The element is more usually produced as a by-product of copper and nickel mining. The Copperbelt in the Democratic Republic of the Congo (DRC) and Zambia yields most of the global cobalt production. World production in 2016 was (according to Natural Resources Canada), and the DRC alone accounted for more than 50%.
Cobalt is primarily used in lithium-ion batteries, and in the manufacture of magnetic, wear-resistant and high-strength alloys. The compounds cobalt silicate and cobalt(II) aluminate (CoAl2O4, cobalt blue) give a distinctive deep blue color to glass, ceramics, inks, paints and varnishes. Cobalt occurs naturally as only one stable isotope, cobalt-59. Cobalt-60 is a commercially important radioisotope, used as a radioactive tracer and for the production of high-energy gamma rays. Cobalt is also used in the petroleum industry as a catalyst when refining crude oil. This is to purge it of sulfur, which is very polluting when burned and causes acid rain.
Cobalt is the active center of a group of coenzymes called cobalamins. Vitamin B, the best-known example of the type, is an essential vitamin for all animals. Cobalt in inorganic form is also a micronutrient for bacteria, algae, and fungi.
The name cobalt derives from a type of ore considered a nuisance by 16th century German silver miners, which in turn may have been named from a spirit or goblin held superstitiously responsible for it; this spirit is considered equitable to the kobold (a household spirit) by some, or, categorized as a gnome (mine spirit) by others.
Characteristics
Cobalt is a ferromagnetic metal with a specific gravity of 8.9. The Curie temperature is and the magnetic moment is 1.6–1.7 Bohr magnetons per atom. Cobalt has a relative permeability two-thirds that of iron. Metallic cobalt occurs as two crystallographic structures: hcp and fcc. The ideal transition temperature between the hcp and fcc structures is , but in practice the energy difference between them is so small that random intergrowth of the two is common.
Cobalt is a weakly reducing metal that is protected from oxidation by a passivating oxide film. It is attacked by halogens and sulfur. Heating in oxygen produces Co3O4 which loses oxygen at to give the monoxide CoO. The metal reacts with fluorine (F2) at 520 K to give CoF3; with chlorine (Cl2), bromine (Br2) and iodine (I2), producing equivalent binary halides. It does not react with hydrogen gas (H2) or nitrogen gas (N2) even when heated, but it does react with boron, carbon, phosphorus, arsenic and sulfur. At ordinary temperatures, it reacts slowly with mineral acids, and very slowly with moist, but not dry, air.
Compounds
Common oxidation states of cobalt include +2 and +3, although compounds with oxidation states ranging from −3 to +5 are also known. A common oxidation state for simple compounds is +2 (cobalt(II)). These salts form the pink-colored metal aquo complex in water. Addition of chloride gives the intensely blue . In a borax bead flame test, cobalt shows deep blue in both oxidizing and reducing flames.
Oxygen and chalcogen compounds
Several oxides of cobalt are known. Green cobalt(II) oxide (CoO) has rocksalt structure. It is readily oxidized with water and oxygen to brown cobalt(III) hydroxide (Co(OH)3). At temperatures of 600–700 °C, CoO oxidizes to the blue cobalt(II,III) oxide (Co3O4), which has a spinel structure. Black cobalt(III) oxide (Co2O3) is also known. Cobalt oxides are antiferromagnetic at low temperature: CoO (Néel temperature 291 K) and Co3O4 (Néel temperature: 40 K), which is analogous to magnetite (Fe3O4), with a mixture of +2 and +3 oxidation states.
The principal chalcogenides of cobalt are the black cobalt(II) sulfides, CoS2 (pyrite structure), (spinel structure), and CoS (nickel arsenide structure).
Halides
Four dihalides of cobalt(II) are known: cobalt(II) fluoride (CoF2, pink), cobalt(II) chloride (CoCl2, blue), cobalt(II) bromide (CoBr2, green), cobalt(II) iodide (CoI2, blue-black). These halides exist in anhydrous and hydrated forms. Whereas the anhydrous dichloride is blue, the hydrate is red.
The reduction potential for the reaction + e− → is +1.92 V, beyond that for chlorine to chloride, +1.36 V. Consequently, cobalt(III) chloride would spontaneously reduce to cobalt(II) chloride and chlorine. Because the reduction potential for fluorine to fluoride is so high, +2.87 V, cobalt(III) fluoride is one of the few simple stable cobalt(III) compounds. Cobalt(III) fluoride, which is used in some fluorination reactions, reacts vigorously with water.
Coordination compounds
The inventory of complexes is very large. Starting with higher oxidation states, complexes of Co(IV) and Co(V) are rare. Examples are found in caesium hexafluorocobaltate(IV) (Cs2CoF6) and potassium percobaltate (K3CoO4).
Cobalt(III) forms a wide variety of coordination complexes with ammonia and amines, which are called ammine complexes. Examples include , (chloropentamminecobalt(III)), and cis- and trans-. The corresponding ethylenediamine complexes are also well known. Analogues are known where the halides are replaced by nitrite, hydroxide, carbonate, etc. Alfred Werner worked extensively on these complexes in his Nobel-prize winning work. The robustness of these complexes is demonstrated by the optical resolution of tris(ethylenediamine)cobalt(III) ().
Cobalt(II) forms a wide variety of complexes, but mainly with weakly basic ligands. The pink-colored cation hexaaquocobalt(II) is found in several routine cobalt salts such as the nitrate and sulfate. Upon addition of excess chloride, solutions of the hexaaquo complex converts to the deep blue , which is tetrahedral.
Softer ligands like triphenylphosphine form complexes with Co(II) and Co(I), examples being bis- and tris(triphenylphosphine)cobalt(I) chloride, and . These Co(I) and Co(II) complexes represent a link to the organometallic complexes described below.
Organometallic compounds
Cobaltocene is a structural analog to ferrocene, with cobalt in place of iron. Cobaltocene is much more sensitive to oxidation than ferrocene. Cobalt carbonyl (Co2(CO)8) is a catalyst in carbonylation and hydrosilylation reactions. Vitamin B12 (see below) is an organometallic compound found in nature and is the only vitamin that contains a metal atom. An example of an alkylcobalt complex in the otherwise uncommon +4 oxidation state of cobalt is the homoleptic complex tetrakis(1-norbornyl)cobalt(IV) (Co(1-norb)4), a transition metal-alkyl complex that is notable for its resistance to β-hydrogen elimination, in accord with Bredt's rule. The cobalt(III) and cobalt(V) complexes and are also known.
Isotopes
59Co is the only stable cobalt isotope and the only isotope that exists naturally on Earth. Twenty-two radioisotopes have been characterized: the most stable, 60Co, has a half-life of 5.2714 years; 57Co has a half-life of 271.8 days; 56Co has a half-life of 77.27 days; and 58Co has a half-life of 70.86 days. All the other radioactive isotopes of cobalt have half-lives shorter than 18 hours, and in most cases shorter than 1 second. This element also has 4 meta states, all of which have half-lives shorter than 15 minutes.
The isotopes of cobalt range in atomic weight from 50 u (50Co) to 73 u (73Co). The primary decay mode for isotopes with atomic mass unit values less than that of the only stable isotope, 59Co, is electron capture and the primary mode of decay in isotopes with atomic mass greater than 59 atomic mass units is beta decay. The primary decay products below 59Co are element 26 (iron) isotopes; above that the decay products are element 28 (nickel) isotopes.
Etymology
Many different stories about the origin of the word "cobalt" have been proposed. In one version the element cobalt was named after "", the name which 16th century German silver miners had given to a nuisance type of ore which occurred that was corrosive and issued poisonous gas. Although such ores had been used for blue pigmentation since antiquity, the Germans at that time did not have the technology to smelt the ore into metal (cf. below).
The authority on such kobelt ore (Latinized as cobaltum or cadmia) at the time was Georgius Agricola. He was also the oft-quoted authority on the mine spirits called "" (Latinized as cobalus or pl. cobali) in a separate work.
Agricola did not make an connection between the similarly named ore and spirit. However, a causal connection (ore blamed on "kobel") was made by a contemporary, and a word origin connection (word "formed" from cobalus) made by a late 18th century writer. Later, Grimms' dictionary (1868) noted the kobalt/kobelt ore was blamed on the mountain spirit () which was also held responsible for "stealing the silver and putting out an ore that caused poor mining atmosphere (Wetter) and other health hazards".
Grimms' dictionary entries equated the word "kobel" with "kobold", and listed it as a mere variant diminutive, but the latter is defined in it as a household spirit. Whereas some of the more recent commentators prefer to characterize the ore's namesake kobelt (recté kobel) as a gnome.
The early 20th century Oxford English Dictionary (1st edition, 1908) had upheld Grimm's etymology. However, by around the same time in Germany, the alternate etymology not endorsed by Grimm (kob/kof "house, chamber" + walt "power, ruler") was being proposed as more convincing.
Somewhat later, Paul Kretschmer (1928) explained that while this "house ruler" etymology was the proper one that backed the original meaning of kobold as household spirit, a corruption later occurred introducing the idea of "mine demon" to it. The present edition of the Etymologisches Wörterbuch (25th ed., 2012) under "kobold" lists the latter, not Grimm's etymology, but still persists, under its entry for "kobalt", that while the cobalt ore may have got its name from "a type of mine spirit/demon" (daemon metallicus) while stating that this is "apparently" the kobold.
Joseph William Mellor (1935) also stated that cobalt may derive from kobalos (), though other theories had been suggested.
Alternate theories
Several alternative etymologies that have been suggested, which may not involve a spirit (kobel or kobold) at all. Karl Müller-Fraureuth conjectured that kobelt derived from , a bucket used in mining, frequently mentioned by Agricola, namely the kobel/köbel (Latinized as modulus).
Another theory given by the Etymologisches Wörterbuch derives the term from or rather (, "arsenic sulfide") which occurs as noxious fumes.
An etymology from Slavonic was suggested by Emanuel Merck (1902).
W. W. Skeat and J. Berendes construed as "parasite", i.e. as an ore parasitic to nickel, but this explanation is faulted for its anachronism since nickel was not discovered until 1751.
History
Cobalt compounds have been used for centuries to impart a rich blue color to glass, glazes, and ceramics. Cobalt has been detected in Egyptian sculpture, Persian jewelry from the third millennium BC, in the ruins of Pompeii, destroyed in 79 AD, and in China, dating from the Tang dynasty (618–907 AD) and the Ming dynasty (1368–1644 AD).
Cobalt has been used to color glass since the Bronze Age. The excavation of the Uluburun shipwreck yielded an ingot of blue glass, cast during the 14th century BC. Blue glass from Egypt was either colored with copper, iron, or cobalt. The oldest cobalt-colored glass is from the eighteenth dynasty of Egypt (1550–1292 BC). The source for the cobalt the Egyptians used is not known.
The word cobalt is derived from the 16th century German "", a type of ore, as aforementioned. The first attempts to smelt those ores for copper or silver failed, yielding simply powder (cobalt(II) oxide) instead. Because the primary ores of cobalt always contain arsenic, smelting the ore oxidized the arsenic into the highly toxic and volatile arsenic oxide, adding to the notoriety of the ore. Paracelsus, Georgius Agricola, and Basil Valentine all referred to such silicates as "cobalt".
Swedish chemist Georg Brandt (1694–1768) is credited with discovering cobalt , showing it to be a previously unknown element, distinct from bismuth and other traditional metals. Brandt called it a new "semi-metal", naming it for the mineral from which he had extracted it. He showed that compounds of cobalt metal were the source of the blue color in glass, which previously had been attributed to the bismuth found with cobalt. Cobalt became the first metal to be discovered since the pre-historical period. All previously known metals (iron, copper, silver, gold, zinc, mercury, tin, lead and bismuth) had no recorded discoverers.
During the 19th century, a significant part of the world's production of cobalt blue (a pigment made with cobalt compounds and alumina) and smalt (cobalt glass powdered for use for pigment purposes in ceramics and painting) was carried out at the Norwegian Blaafarveværket. The first mines for the production of smalt in the 16th century were located in Norway, Sweden, Saxony and Hungary. With the discovery of cobalt ore in New Caledonia in 1864, the mining of cobalt in Europe declined. With the discovery of ore deposits in Ontario, Canada, in 1904 and the discovery of even larger deposits in the Katanga Province in the Congo in 1914, mining operations shifted again. When the Shaba conflict started in 1978, the copper mines of Katanga Province nearly stopped production. The impact on the world cobalt economy from this conflict was smaller than expected: cobalt is a rare metal, the pigment is highly toxic, and the industry had already established effective ways for recycling cobalt materials. In some cases, industry was able to change to cobalt-free alternatives.
In 1938, John Livingood and Glenn T. Seaborg discovered the radioisotope cobalt-60. This isotope was famously used at Columbia University in the 1950s to establish parity violation in radioactive beta decay.
After World War II, the US wanted to guarantee the supply of cobalt ore for military uses (as the Germans had been doing) and prospected for cobalt within the US. High purity cobalt was highly sought after for its use in jet engines and gas turbines. An adequate supply of the ore was found in Idaho near Blackbird canyon. Calera Mining Company started production at the site.
Cobalt demand has further accelerated in the 21st century as an essential constituent of materials used in rechargeable batteries, superalloys, and catalysts. It has been argued that cobalt will be one of the main objects of geopolitical competition in a world running on renewable energy and dependent on batteries, but this perspective has also been criticised for underestimating the power of economic incentives for expanded production.
Occurrence
The stable form of cobalt is produced in supernovae through the r-process. It comprises 0.0029% of the Earth's crust. Except as recently delivered in meteoric iron, free cobalt (the native metal) is not found on Earth's surface because of its tendency to react with oxygen in the atmosphere. Small amounts of cobalt compounds are found in most rocks, soils, plants, and animals. In the ocean cobalt typically reacts with chlorine.
In nature, cobalt is frequently associated with nickel. Both are characteristic components of meteoric iron, though cobalt is much less abundant in iron meteorites than nickel. As with nickel, cobalt in meteoric iron alloys may have been well enough protected from oxygen and moisture to remain as the free (but alloyed) metal.
Cobalt in compound form occurs in copper and nickel minerals. It is the major metallic component that combines with sulfur and arsenic in the sulfidic cobaltite (CoAsS), safflorite (CoAs2), glaucodot (), and skutterudite (CoAs3) minerals. The mineral cattierite is similar to pyrite and occurs together with vaesite in the copper deposits of Katanga Province. When it reaches the atmosphere, weathering occurs; the sulfide minerals oxidize and form pink erythrite ("cobalt glance": Co3(AsO4)2·8H2O) and spherocobaltite (CoCO3).
Cobalt is also a constituent of tobacco smoke. The tobacco plant readily absorbs and accumulates heavy metals like cobalt from the surrounding soil in its leaves. These are subsequently inhaled during tobacco smoking.
Production
The main ores of cobalt are cobaltite, erythrite, glaucodot and skutterudite (see above), but most cobalt is obtained by reducing the cobalt by-products of nickel and copper mining and smelting.
Since cobalt is generally produced as a by-product, the supply of cobalt depends to a great extent on the economic feasibility of copper and nickel mining in a given market. Demand for cobalt was projected to grow 6% in 2017.
Primary cobalt deposits are rare, such as those occurring in hydrothermal deposits, associated with ultramafic rocks, typified by the Bou-Azzer district of Morocco. At such locations, cobalt ores are mined exclusively, albeit at a lower concentration, and thus require more downstream processing for cobalt extraction.
Several methods exist to separate cobalt from copper and nickel, depending on the concentration of cobalt and the exact composition of the used ore. One method is froth flotation, in which surfactants bind to ore components, leading to an enrichment of cobalt ores. Subsequent roasting converts the ores to cobalt sulfate, and the copper and the iron are oxidized to the oxide. Leaching with water extracts the sulfate together with the arsenates. The residues are further leached with sulfuric acid, yielding a solution of copper sulfate. Cobalt can also be leached from the slag of copper smelting.
The products of the above-mentioned processes are transformed into the cobalt oxide (Co3O4). This oxide is reduced to metal by the aluminothermic reaction or reduction with carbon in a blast furnace.
Extraction
The United States Geological Survey estimates world reserves of cobalt at 7,100,000 metric tons. The Democratic Republic of the Congo (DRC) currently produces 63% of the world's cobalt. This market share may reach 73% by 2025 if planned expansions by mining producers like Glencore Plc take place as expected. Bloomberg New Energy Finance has estimated that by 2030, global demand for cobalt could be 47 times more than it was in 2017.
Democratic Republic of the Congo
Changes that Congo made to mining laws in 2002 attracted new investments in Congolese copper and cobalt projects. In 2005, the top producer of cobalt was the copper deposits in the Democratic Republic of the Congo's Katanga Province. Formerly Shaba province, the area had almost 40% of global reserves, reported the British Geological Survey in 2009.
The Mukondo Mountain project, operated by the Central African Mining and Exploration Company (CAMEC) in Katanga Province, may be the richest cobalt reserve in the world. It produced an estimated one-third of the total global cobalt production in 2008. In July 2009, CAMEC announced a long-term agreement to deliver its entire annual production of cobalt concentrate from Mukondo Mountain to Zhejiang Galico Cobalt & Nickel Materials of China.
In 2016, Chinese ownership of cobalt production in the Congo was estimated at over 10% of global cobalt supply, forming a key input to the Chinese cobalt refining industry and granting China substantial influence over the global cobalt supply chain. Chinese control of Congolese cobalt has raised concern in Western nations which have sought to reduce supply chain reliance upon China and have expressed concern regarding labor and human rights violations in cobalt mines in the DRC.
Glencore's Mutanda Mine shipped 24,500 tons of cobalt in 2016, 40% of Congo DRC's output and nearly a quarter of global production. After oversupply, Glencore closed Mutanda for two years in late 2019. Glencore's Katanga Mining project is resuming as well and should produce 300,000 tons of copper and 20,000 tons of cobalt by 2019, according to Glencore.
In February 2018, global asset management firm AllianceBernstein defined the DRC as economically "the Saudi Arabia of the electric vehicle age", due to its cobalt resources, as essential to the lithium-ion batteries that drive electric vehicles.
On March 9, 2018, President Joseph Kabila updated the 2002 mining code, increasing royalty charges and declaring cobalt and coltan "strategic metals". The 2002 mining code was effectively updated on December 4, 2018.
Labor conditions
Artisanal mining supplied 17% to 40% of the DRC production as of 2016. Some 100,000 cobalt miners in Congo DRC use hand tools to dig hundreds of feet, with little planning and fewer safety measures, say workers and government and NGO officials, as well as The Washington Post reporters' observations on visits to isolated mines. The lack of safety precautions frequently causes injuries or death. Mining pollutes the vicinity and exposes local wildlife and indigenous communities to toxic metals thought to cause birth defects and breathing difficulties, according to health officials.
Child labor is used in mining cobalt from African artisanal mines. Human rights activists have highlighted this and investigative journalism reporting has confirmed it. This revelation prompted cell phone maker Apple Inc., on March 3, 2017, to stop buying ore from suppliers such as Zhejiang Huayou Cobalt who source from artisanal mines in the DRC, and begin using only suppliers that are verified to meet its workplace standards. In 2023, Apple announced it would convert to using recycled cobalt by 2025.
There is a push globally by the EU and major car manufacturers (OEM) for global production of cobalt to be sourced and –produced sustainably, responsibly and traceability of the supply chain. Mining companies are adopting and practising ESG initiatives in line with OECD Guidance and putting in place evidence of zero to low carbon footprint activities in the supply chain production of lithium-ion batteries. These initiatives are already taking place with major mining companies, artisanal and small-scale mining companies (ASM). Car manufacturers and battery manufacturer supply chains: Tesla, VW, BMW, BASF and Glencore are participating in several initiatives, such as the Responsible Cobalt Initiative and Cobalt for Development study. In 2018 BMW Group in partnership with BASF, Samsung SDI and Samsung Electronics have launched a pilot project in the DRC over one pilot mine, to improve conditions and address challenges for artisanal miners and the surrounding communities.
The political and ethnic dynamics of the region have in the past caused outbreaks of violence and years of armed conflict and displaced populations. This instability affected the price of cobalt and also created perverse incentives for the combatants in the First and Second Congo Wars to prolong the fighting, since access to diamond mines and other valuable resources helped to finance their military goals—which frequently amounted to genocide—and also enriched the fighters themselves. While DR Congo has in the 2010s not recently been invaded by neighboring military forces, some of the richest mineral deposits adjoin areas where Tutsis and Hutus still frequently clash, unrest continues although on a smaller scale and refugees still flee outbreaks of violence.
Cobalt extracted from small Congolese artisanal mining endeavors in 2007 supplied a single Chinese company, Congo DongFang International Mining. A subsidiary of Zhejiang Huayou Cobalt, one of the world's largest cobalt producers, Congo DongFang supplied cobalt to some of the world's largest battery manufacturers, who produced batteries for ubiquitous products like the Apple iPhones. Because of accused labour violations and environmental concerns, LG Chem subsequently audited Congo DongFang in accordance with OECD guidelines. LG Chem, which also produces battery materials for car companies, imposed a code of conduct on all suppliers that it inspects.
In December 2019, International Rights Advocates, a human rights NGO, filed a landmark lawsuit against Apple, Tesla, Dell, Microsoft and Google company Alphabet for "knowingly benefiting from and aiding and abetting the cruel and brutal use of young children" in mining cobalt. The companies in question denied their involvement in child labour. In 2024 the court ruled that the suppliers facilitate force labor but the US tech companies are not liable because they don't operate as a shared enterprise with the suppliers and that the "alleged injuries are not fairly traceable" to any of the defendants' conduct. The book Cobalt Red alleges that workers including children suffer injuries, amputations, and death as the result of the hazardous working conditions and mine tunnel collapses during artisanal mining of cobalt in the DRC.
Since child and slave labor have been repeatedly reported in cobalt mining, primarily in the artisanal mines of DR Congo, technology companies seeking an ethical supply chain have faced shortages of this raw material and the price of cobalt metal reached a nine-year high in October 2017, more than US$30 a pound, versus US$10 in late 2015. After oversupply, the price dropped to a more normal $15 in 2019. As a reaction to the issues with artisanal cobalt mining in DR Congo a number of cobalt suppliers and their customers have formed the Fair Cobalt Alliance (FCA) which aims to end the use of child labor and to improve the working conditions of cobalt mining and processing in the DR Congo. Members of FCA include Zhejiang Huayou Cobalt, Sono Motors, the Responsible Cobalt Initiative, Fairphone, Glencore and Tesla, Inc.
Canada
In 2017, some exploration companies were planning to survey old silver and cobalt mines in the area of Cobalt, Ontario, where significant deposits are believed to lie.
Cuba
Canada's Sherritt International processes cobalt ores in nickel deposits from the Moa mines in Cuba, and the island has several others mines in Mayarí, Camagüey, and Pinar del Río. Continued investments by Sherritt International in Cuban nickel and cobalt production while acquiring mining rights for 17–20 years made the communist country third for cobalt reserves in 2019, before Canada itself.
Indonesia
Starting from smaller amounts in 2021, Indonesia began producing cobalt as a byproduct of nickel production. By 2022, the country had become the world's second-largest cobalt producer, with Benchmark Mineral Intelligence forecasting Indonesian output to make up 20 percent of global production by 2030.
Applications
In 2016, of cobalt was used. Cobalt has been used in the production of high-performance alloys. It is also used in some rechargeable batteries.
Alloys
Cobalt-based superalloys have historically consumed most of the cobalt produced. The temperature stability of these alloys makes them suitable for turbine blades for gas turbines and aircraft jet engines, although nickel-based single-crystal alloys surpass them in performance. Cobalt-based alloys are also corrosion- and wear-resistant, making them, like titanium, useful for making orthopedic implants that do not wear down over time. The development of wear-resistant cobalt alloys started in the first decade of the 20th century with the stellite alloys, containing chromium with varying quantities of tungsten and carbon. Alloys with chromium and tungsten carbides are very hard and wear-resistant. Special cobalt-chromium-molybdenum alloys like Vitallium are used for prosthetic parts (hip and knee replacements). Cobalt alloys are also used for dental prosthetics as a useful substitute for nickel, which may be allergenic. Some high-speed steels also contain cobalt for increased heat and wear resistance. The special alloys of aluminium, nickel, cobalt and iron, known as Alnico, and of samarium and cobalt (samarium–cobalt magnet) are used in permanent magnets. It is also alloyed with 95% platinum for jewelry, yielding an alloy suitable for fine casting, which is also slightly magnetic.
Batteries
Lithium cobalt oxide (LiCoO2, aka "LCO"), first sold commercially in 1991 by Sony, was widely used in lithium-ion battery cathodes until the 2010s. The material is composed of cobalt oxide layers with the lithium intercalated. These LCO batteries continue to dominate the market for consumer electronics. Batteries for electric cars however have shifted to lower cobalt technologies.
In 2018 most cobalt in batteries was used in a mobile device, a more recent application for cobalt is rechargeable batteries for electric cars. This industry increased five-fold in its demand for cobalt from 2016 to 2020, which made it urgent to find new raw materials in more stable areas of the world. Demand is expected to continue or increase as the prevalence of electric vehicles increases. Exploration in 2016–2017 included the area around Cobalt, Ontario, an area where many silver mines ceased operation decades ago. Cobalt for electric vehicles increased 81% from the first half of 2018 to 7,200 tonnes in the first half of 2019, for a battery capacity of 46.3 GWh.
As of August 2020 battery makers have gradually reduced the cathode cobalt content from 1/3 (NMC 111) to 1/5 (NMC 442) to currently 1/10 (NMC 811) and have also introduced the cobalt free lithium iron phosphate cathode into the battery packs of electric cars such as the Tesla Model 3.
Research was also conducted by the European Union into the possibility of eliminating cobalt requirements in lithium-ion battery production.
In September 2020, Tesla outlined their plans to make their own, cobalt-free battery cells.
Nickel–cadmium (NiCd) and nickel metal hydride (NiMH) batteries also included cobalt to improve the oxidation of nickel in the battery.
Lithium iron phosphate batteries officially surpassed ternary cobalt batteries in 2021 with 52% of installed capacity. Analysts estimate that its market share will exceed 60% in 2024.
Catalysts
Several cobalt compounds are oxidation catalysts. Cobalt acetate is used to convert xylene to terephthalic acid, the precursor of the bulk polymer polyethylene terephthalate. Typical catalysts are the cobalt carboxylates (known as cobalt soaps). They are also used in paints, varnishes, and inks as "drying agents" through the oxidation of drying oils. However, their use is being phased out due to toxicity concerns. The same carboxylates are used to improve the adhesion between steel and rubber in steel-belted radial tires. In addition they are used as accelerators in polyester resin systems.
Cobalt-based catalysts are used in reactions involving carbon monoxide. Cobalt is also a catalyst in the Fischer–Tropsch process for the hydrogenation of carbon monoxide into liquid fuels. Hydroformylation of alkenes often uses cobalt octacarbonyl as a catalyst. The hydrodesulfurization of petroleum uses a catalyst derived from cobalt and molybdenum. This process helps to clean petroleum of sulfur impurities that interfere with the refining of liquid fuels.
Pigments and coloring
Before the 19th century, cobalt was predominantly used as a pigment. It has been used since the Middle Ages to make smalt, a blue-colored glass. Smalt is produced by melting a mixture of roasted mineral smaltite, quartz and potassium carbonate, which yields a dark blue silicate glass, which is finely ground after the production. Smalt was widely used to color glass and as pigment for paintings. In 1780, Sven Rinman discovered cobalt green, and in 1802 Louis Jacques Thénard discovered cobalt blue. Cobalt pigments such as cobalt blue (cobalt aluminate), cerulean blue (cobalt(II) stannate), various hues of cobalt green (a mixture of cobalt(II) oxide and zinc oxide), and cobalt violet (cobalt phosphate) are used as artist's pigments because of their superior chromatic stability.
Radioisotopes
Cobalt-60 (Co-60 or 60Co) is useful as a gamma-ray source because it can be produced in predictable amounts with high activity by bombarding cobalt with neutrons. It produces gamma rays with energies of 1.17 and 1.33 MeV.
Cobalt is used in external beam radiotherapy, sterilization of medical supplies and medical waste, radiation treatment of foods for sterilization (cold pasteurization), industrial radiography (e.g. weld integrity radiographs), density measurements (e.g. concrete density measurements), and tank fill height switches. The metal has the unfortunate property of producing a fine dust, causing problems with radiation protection. Cobalt from radiotherapy machines has been a serious hazard when not discarded properly, and one of the worst radiation contamination accidents in North America occurred in 1984, when a discarded radiotherapy unit containing cobalt-60 was mistakenly disassembled in a junkyard in Juarez, Mexico.
Cobalt-60 has a radioactive half-life of 5.27 years. Loss of potency requires periodic replacement of the source in radiotherapy and is one reason why cobalt machines have been largely replaced by linear accelerators in modern radiation therapy. Cobalt-57 (Co-57 or 57Co) is a cobalt radioisotope most often used in medical tests, as a radiolabel for vitamin B uptake, and for the Schilling test. Cobalt-57 is used as a source in Mössbauer spectroscopy and is one of several possible sources in X-ray fluorescence devices.
Nuclear weapon designs could intentionally incorporate 59Co, some of which would be activated in a nuclear explosion to produce 60Co. The 60Co, dispersed as nuclear fallout, is sometimes called a cobalt bomb.
Magnetic materials
Due to the ferromagnetic properties of cobalt, it is used in the production of various magnetic materials. It is used in creating permanent magnets like Alnico magnets, known for their strong magnetic properties used in electric motors, sensors, and MRI machines. It is also used in production of magnetic alloys like cobalt steel, widely used in magnetic recording media such as hard disks and tapes.
Cobalt's ability to maintain magnetic properties at high temperatures makes it valuable in magnetic recording applications, ensuring reliable data storage devices. Cobalt also contributes to specialized magnets such as samarium-cobalt and neodymium-iron-boron magnets, which are vital in electronics for components like sensors and actuators.
Other uses
Cobalt is used in electroplating for its attractive appearance, hardness, and resistance to oxidation.
It is also used as a base primer coat for porcelain enamels.
Biological role
Cobalt is essential to the metabolism of all animals. It is a key constituent of cobalamin, also known as vitamin B, the primary biological reservoir of cobalt as an ultratrace element. Bacteria in the stomachs of ruminant animals convert cobalt salts into vitamin B, a compound which can only be produced by bacteria or archaea. A minimal presence of cobalt in soils therefore markedly improves the health of grazing animals, and an uptake of 0.20 mg/kg a day is recommended, because they have no other source of vitamin B.
Proteins based on cobalamin use corrin to hold the cobalt. Coenzyme B12 features a reactive C-Co bond that participates in the reactions. In humans, B12 has two types of alkyl ligand: methyl and adenosyl. MeB12 promotes methyl (−CH3) group transfers. The adenosyl version of B12 catalyzes rearrangements in which a hydrogen atom is directly transferred between two adjacent atoms with concomitant exchange of the second substituent, X, which may be a carbon atom with substituents, an oxygen atom of an alcohol, or an amine. Methylmalonyl coenzyme A mutase (MUT) converts MMl-CoA to Su-CoA, an important step in the extraction of energy from proteins and fats.
Although far less common than other metalloproteins (e.g. those of zinc and iron), other cobaltoproteins are known besides B12. These proteins include methionine aminopeptidase 2, an enzyme that occurs in humans and other mammals that does not use the corrin ring of B12, but binds cobalt directly. Another non-corrin cobalt enzyme is nitrile hydratase, an enzyme in bacteria that metabolizes nitriles.
Cobalt deficiency
In humans, consumption of cobalt-containing vitamin B12 meets all needs for cobalt. For cattle and sheep, which meet vitamin B12 needs via synthesis by resident bacteria in the rumen, there is a function for inorganic cobalt. In the early 20th century, during the development of farming on the North Island Volcanic Plateau of New Zealand, cattle suffered from what was termed "bush sickness". It was discovered that the volcanic soils lacked the cobalt salts essential for the cattle food chain. The "coast disease" of sheep in the Ninety Mile Desert of the Southeast of South Australia in the 1930s was found to originate in nutritional deficiencies of trace elements cobalt and copper. The cobalt deficiency was overcome by the development of "cobalt bullets", dense pellets of cobalt oxide mixed with clay given orally for lodging in the animal's rumen.
Health issues
The LD50 value for soluble cobalt salts has been estimated to be between 150 and 500 mg/kg. In the US, the Occupational Safety and Health Administration (OSHA) has designated a permissible exposure limit (PEL) in the workplace as a time-weighted average (TWA) of 0.1 mg/m3. The National Institute for Occupational Safety and Health (NIOSH) has set a recommended exposure limit (REL) of 0.05 mg/m3, time-weighted average. The IDLH (immediately dangerous to life and health) value is 20 mg/m3.
However, chronic cobalt ingestion has caused serious health problems at doses far less than the lethal dose. In 1966, the addition of cobalt compounds to stabilize beer foam in Canada led to a peculiar form of toxin-induced cardiomyopathy, which came to be known as beer drinker's cardiomyopathy.
Furthermore, cobalt metal is suspected of causing cancer (i.e., possibly carcinogenic, IARC Group 2B) as per the International Agency for Research on Cancer (IARC) Monographs.
It causes respiratory problems when inhaled. It also causes skin problems when touched; after nickel and chromium, cobalt is a major cause of contact dermatitis.
| Physical sciences | Chemical elements_2 | null |
24580596 | https://en.wikipedia.org/wiki/Cerium | Cerium | Cerium is a chemical element; it has symbol Ce and atomic number 58. It is a soft, ductile, and silvery-white metal that tarnishes when exposed to air. Cerium is the second element in the lanthanide series, and while it often shows the oxidation state of +3 characteristic of the series, it also has a stable +4 state that does not oxidize water. It is considered one of the rare-earth elements. Cerium has no known biological role in humans but is not particularly toxic, except with intense or continued exposure.
Despite always occurring in combination with the other rare-earth elements in minerals such as those of the monazite and bastnäsite groups, cerium is easy to extract from its ores, as it can be distinguished among the lanthanides by its unique ability to be oxidized to the +4 state in aqueous solution. It is the most common of the lanthanides, followed by neodymium, lanthanum, and praseodymium. Its estimated abundance in the Earth's crust is 68 ppm.
Cerium was the first of the lanthanides to be discovered, in Bastnäs, Sweden. It was discovered by Jöns Jakob Berzelius and Wilhelm Hisinger in 1803, and independently by Martin Heinrich Klaproth in Germany in the same year. In 1839 Carl Gustaf Mosander became the first to isolate the metal. Today, cerium and its compounds have a variety of uses: for example, cerium(IV) oxide is used to polish glass and is an important part of catalytic converters. Cerium metal is used in ferrocerium lighters for its pyrophoric properties. Cerium-doped YAG phosphor is used in conjunction with blue light-emitting diodes to produce white light in most commercial white LED light sources.
Characteristics
Physical
Cerium is the second element of the lanthanide series. In the periodic table, it appears between the lanthanides lanthanum to its left and praseodymium to its right, and above the actinide thorium. It is a ductile metal with a hardness similar to that of silver. Its 58 electrons are arranged in the configuration [Xe]4f5d6s, of which the four outer electrons are valence electrons. The 4f, 5d, and 6s energy levels are very close to each other, and the transfer of one electron to the 5d shell is due to strong interelectronic repulsion in the compact 4f shell. This effect is overwhelmed when the atom is positively ionised; thus Ce on its own has instead the regular configuration [Xe]4f, although in some solid solutions it may be [Xe]4f5d. Most lanthanides can use only three electrons as valence electrons, as afterwards the remaining 4f electrons are too strongly bound: cerium is an exception because of the stability of the empty f-shell in Ce and the fact that it comes very early in the lanthanide series, where the nuclear charge is still low enough until neodymium to allow the removal of the fourth valence electron by chemical means.
Cerium has a variable electronic structure. The energy of the 4f electron is nearly the same as that of the outer 5d and 6s electrons that are delocalized in the metallic state, and only a small amount of energy is required to change the relative occupancy of these electronic levels. This gives rise to dual valence states. For example, a volume change of about 10% occurs when cerium is subjected to high pressures or low temperatures. In its high pressure phase (α-Cerium), the 4f electrons are also delocalized and itinerate, as opposed to localized 4f electrons in low pressure phase (γ-Cerium). It appears that the valence changes from about 3 to 4 when it is cooled or compressed.
Chemical properties of the element
Like the other lanthanides, cerium metal is a good reducing agent, having standard reduction potential of E = −2.34 V for the Ce/Ce couple. It tarnishes in air, forming a passivating oxide layer like iron rust. A centimeter-sized sample of cerium metal corrodes completely in about a year. More dramatically, metallic cerium can be highly pyrophoric:
Being highly electropositive, cerium reacts with water. The reaction is slow with cold water but speeds up with increasing temperature, producing cerium(III) hydroxide and hydrogen gas:
Allotropes
Four allotropic forms of cerium are known to exist at standard pressure and are given the common labels of α to δ:
The high-temperature form, δ-cerium, has a bcc (body-centered cubic) crystal structure and exists above 726 °C.
The stable form below 726 °C to approximately room temperature is γ-cerium, with an fcc (face-centered cubic) crystal structure.
The DHCP (double hexagonal close-packed) form β-cerium is the equilibrium structure approximately from room temperature to −150 °C.
The fcc form α-cerium is stable below about −150 °C; it has a density of 8.16 g/cm.
Other solid phases occurring only at high pressures are shown on the phase diagram.
Both γ and β forms are quite stable at room temperature, although the equilibrium transformation temperature is estimated at 75 °C.
At lower temperatures the behavior of cerium is complicated by the slow rates of transformation. Transformation temperatures are subject to substantial hysteresis and values quoted here are approximate. Upon cooling below −15 °C, γ-cerium starts to change to β-cerium, but the transformation involves a volume increase and, as more β forms, the internal stresses build up and suppress further transformation. Cooling below approximately −160 °C will start formation of α-cerium but this is only from remaining γ-cerium. β-cerium does not significantly transform to α-cerium except in the presence of stress or deformation. At atmospheric pressure, liquid cerium is more dense than its solid form at the melting
point.
Isotopes
Naturally occurring cerium is made up of four isotopes: Ce (0.19%), Ce (0.25%), Ce (88.4%), and Ce (11.1%). All four are observationally stable, though the light isotopes Ce and Ce are theoretically expected to undergo double electron capture to isotopes of barium, and the heaviest isotope 142Ce is expected to undergo double beta decay to 142Nd or alpha decay to 138Ba. Thus, 140Ce is the only theoretically stable isotope. None of these decay modes have yet been observed, though the double beta decay of 136Ce, 138Ce, and 142Ce have been experimentally searched for. The current experimental limits for their half-lives are:
136Ce: >3.8×1016 y
138Ce: >5.7×1016 y
142Ce: >5.0×1016 y
All other cerium isotopes are synthetic and radioactive. The most stable of them are 144Ce with a half-life of 284.9 days, 139Ce with a half-life of 137.6 days, and 141Ce with a half-life of 32.5 days. All other radioactive cerium isotopes have half-lives under four days, and most of them have half-lives under ten minutes. The isotopes between 140Ce and 144Ce inclusive occur as fission products of uranium. The primary decay mode of the isotopes lighter than 140Ce is inverse beta decay or electron capture to isotopes of lanthanum, while that of the heavier isotopes is beta decay to isotopes of praseodymium. Some isotopes of neodymium can alpha decay or are predicted to decay to isotopes of cerium.
The rarity of the proton-rich 136Ce and 138Ce is explained by the fact that they cannot be made in the most common processes of stellar nucleosynthesis for elements beyond iron, the s-process (slow neutron capture) and the r-process (rapid neutron capture). This is so because they are bypassed by the reaction flow of the s-process, and the r-process nuclides are blocked from decaying to them by more neutron-rich stable nuclides. Such nuclei are called p-nuclei, and their origin is not yet well understood: some speculated mechanisms for their formation include proton capture as well as photodisintegration. 140Ce is the most common isotope of cerium, as it can be produced in both the s- and r-processes, while 142Ce can only be produced in the r-process. Another reason for the abundance of 140Ce is that it is a magic nucleus, having a closed neutron shell (it has 82 neutrons), and hence it has a very low cross section towards further neutron capture. Although its proton number of 58 is not magic, it is granted additional stability, as its eight additional protons past the magic number 50 enter and complete the 1g proton orbital. The abundances of the cerium isotopes may differ very slightly in natural sources, because Ce and Ce are the daughters of the long-lived primordial radionuclides La and Nd, respectively.
Compounds
Cerium exists in two main oxidation states, Ce(III) and Ce(IV). This pair of adjacent oxidation states dominates several aspects of the chemistry of this element. Cerium(IV) aqueous solutions may be prepared by reacting cerium(III) solutions with the strong oxidizing agents peroxodisulfate or bismuthate. The value of E(Ce/Ce) varies widely depending on conditions due to the relative ease of complexation and hydrolysis with various anions, although +1.72 V is representative. Cerium is the only lanthanide which has important aqueous and coordination chemistry in the +4 oxidation state.
Halides
Cerium forms all four trihalides CeX (X = F, Cl, Br, I) usually by reaction of the oxides with the hydrogen halides. The anhydrous halides are pale-colored, paramagnetic, hygroscopic solids. Upon hydration, the trihalides convert to complexes containing aquo complexes [Ce(HO)]. Unlike most lanthanides, Ce forms a tetrafluoride, a white solid. It also forms a bronze-colored diiodide, which has metallic properties. Aside from the binary halide phases, a number of anionic halide complexes are known. The fluoride gives the Ce(IV) derivatives and . The chloride gives the orange .
Oxides and chalcogenides
Cerium(IV) oxide ("ceria") has the fluorite structure, similarly to the dioxides of praseodymium and terbium. Ceria is a nonstoichiometric compound, meaning that the real formula is CeO, where x is about 0.2. Thus, the material is not perfectly described as Ce(IV). Ceria reduces to cerium(III) oxide with hydrogen gas. Many nonstoichiometric chalcogenides are also known, along with the trivalent CeZ (Z = S, Se, Te). The monochalcogenides CeZ conduct electricity and would better be formulated as CeZe. While CeZ are known, they are polychalcogenides with cerium(III): cerium(IV) derivatives of S, Se, and Te are unknown.
Cerium(IV) complexes
The compound ceric ammonium nitrate (CAN) is the most common cerium compound encountered in the laboratory. The six nitrate ligands bind as bidentate ligands. The complex is 12-coordinate, a high coordination number which emphasizes the large size of the Ce4+ ion. CAN is a popular oxidant in organic synthesis, both as a stoichiometric reagent and as a catalyst. It is inexpensive, stable in air, easily handled, and of low toxicity. It operates by one-electron redox. Cerium nitrates also form 4:3 and 1:1 complexes with 18-crown-6 (the ratio referring to that between the nitrate and the crown ether). Classically, CAN is a primary standard for quantitative analysis. Cerium(IV) salts, especially cerium(IV) sulfate, are often used as standard reagents for volumetric analysis in cerimetric titrations.
Due to ligand-to-metal charge transfer, aqueous cerium(IV) ions are orange-yellow. Aqueous cerium(IV) is metastable in water and is a strong oxidizing agent that oxidizes hydrochloric acid to give chlorine gas. In the Belousov–Zhabotinsky reaction, cerium oscillates between the +4 and +3 oxidation states to catalyze the reaction.
Organocerium compounds
Organocerium chemistry is similar to that of the other lanthanides, often involving complexes of cyclopentadienyl and cyclooctatetraenyl ligands. Cerocene adopts the uranocene molecular structure. The 4f electron in cerocene is poised ambiguously between being localized and delocalized and this compound is considered intermediate-valent.
Alkyl, alkynyl, and alkenyl organocerium derivatives are prepared from the transmetallation of the respective organolithium or Grignard reagents, and are more nucleophilic but less basic than their precursors.
History
Cerium was discovered in Bastnäs in Sweden by Jöns Jakob Berzelius and Wilhelm Hisinger, and independently in Germany by Martin Heinrich Klaproth, both in 1803. Cerium was named by Berzelius after the asteroid Ceres, formally 1 Ceres, discovered two years earlier. Ceres was initially considered to be a planet at the time. The asteroid is itself named after the Roman goddess Ceres, goddess of agriculture, grain crops, fertility and motherly relationships.
Cerium was originally isolated in the form of its oxide, which was named ceria, a term that is still used. The metal itself was too electropositive to be isolated by then-current smelting technology, a characteristic of rare-earth metals in general. After the development of electrochemistry by Humphry Davy five years later, the earths soon yielded the metals they contained. Ceria, as isolated in 1803, contained all of the lanthanides present in the cerite ore from Bastnäs, Sweden, and thus only contained about 45% of what is now known to be pure ceria. It was not until Carl Gustaf Mosander succeeded in removing lanthana and "didymia" in the late 1830s that ceria was obtained pure. Wilhelm Hisinger was a wealthy mine-owner and amateur scientist, and sponsor of Berzelius. He owned and controlled the mine at Bastnäs, and had been trying for years to find out the composition of the abundant heavy gangue rock (the "Tungsten of Bastnäs", which despite its name contained no tungsten), now known as cerite, that he had in his mine. Mosander and his family lived for many years in the same house as Berzelius, and Mosander was undoubtedly persuaded by Berzelius to investigate ceria further.
The element played a role in the Manhattan Project, where cerium compounds were investigated in the Berkeley site as materials for crucibles for uranium and plutonium casting. For this reason, new methods for the preparation and casting of cerium were developed within the scope of the Ames daughter project (now the Ames Laboratory). Production of extremely pure cerium in Ames commenced in mid-1944 and continued until August 1945.
Occurrence and production
Cerium is the most abundant of all the lanthanides and the 25th most abundant element, making up 68 ppm of the Earth's crust. This value is the same of copper, and cerium is even more abundant than common metals such as lead (13 ppm) and tin (2.1 ppm). Thus, despite its position as one of the so-called rare-earth metals, cerium is actually not rare at all. Cerium content in the soil varies between 2 and 150 ppm, with an average of 50 ppm; seawater contains 1.5 parts per trillion of cerium. Cerium occurs in various minerals, but the most important commercial sources are the minerals of the monazite and bastnäsite groups, where it makes up about half of the lanthanide content. Monazite-(Ce) is the most common representative of the monazites, with "-Ce" being the Levinson suffix informing on the dominance of the particular REE element representative. Also the cerium-dominant bastnäsite-(Ce) is the most important of the bastnäsites. Cerium is the easiest lanthanide to extract from its minerals because it is the only one that can reach a stable +4 oxidation state in aqueous solution. Because of the decreased solubility of cerium in the +4 oxidation state, cerium is sometimes depleted from rocks relative to the other rare-earth elements and is incorporated into zircon, since Ce4+ and Zr4+ have the same charge and similar ionic radii. In extreme cases, cerium(IV) can form its own minerals separated from the other rare-earth elements, such as cerianite-(Ce) and .
Bastnäsite, LnCOF, is usually lacking in thorium and the heavy lanthanides beyond samarium and europium, and hence the extraction of cerium from it is quite direct. First, the bastnäsite is purified, using dilute hydrochloric acid to remove calcium carbonate impurities. The ore is then roasted in the air to oxidize it to the lanthanide oxides: while most of the lanthanides will be oxidized to the sesquioxides , cerium will be oxidized to the dioxide CeO. This is insoluble in water and can be leached out with 0.5 M hydrochloric acid, leaving the other lanthanides behind.
The procedure for monazite, , which usually contains all the rare earths, as well as thorium, is more involved. Monazite, because of its magnetic properties, can be separated by repeated electromagnetic separation. After separation, it is treated with hot concentrated sulfuric acid to produce water-soluble sulfates of rare earths. The acidic filtrates are partially neutralized with sodium hydroxide to pH 3–4. Thorium precipitates out of solution as hydroxide and is removed. After that, the solution is treated with ammonium oxalate to convert rare earths to their insoluble oxalates. The oxalates are converted to oxides by annealing. The oxides are dissolved in nitric acid, but cerium oxide is insoluble in HNO3 and hence precipitates out. Care must be taken when handling some of the residues as they contain 228Ra, the daughter of 232Th, which is a strong gamma emitter.
Applications
Cerium has two main applications, both of which use CeO2. The industrial application of ceria is for polishing, especially chemical-mechanical planarization (CMP). In its other main application, CeO2 is used to decolorize glass. It functions by converting green-tinted ferrous impurities to nearly colorless ferric oxides. Ceria has also been used as a substitute for its radioactive congener thoria, for example in the manufacture of electrodes used in gas tungsten arc welding, where ceria as an alloying element improves arc stability and ease of starting while decreasing burn-off.
Gas mantles and pyrophoric alloys
The first use of cerium was in gas mantles, invented by Austrian chemist Carl Auer von Welsbach. In 1885, he had previously experimented with mixtures of magnesium, lanthanum, and yttrium oxides, but these gave green-tinted light and were unsuccessful. Six years later, he discovered that pure thorium oxide produced a much better, though blue, light, and that mixing it with cerium dioxide resulted in a bright white light. Cerium dioxide also acts as a catalyst for the combustion of thorium oxide.
This resulted in commercial success for von Welsbach and his invention, and created great demand for thorium. Its production resulted in a large amount of lanthanides being simultaneously extracted as by-products. Applications were soon found for them, especially in the pyrophoric alloy known as "mischmetal" composed of 50% cerium, 25% lanthanum, and the remainder being the other lanthanides, that is used widely for lighter flints. Usually iron is added to form the alloy ferrocerium, also invented by von Welsbach. Due to the chemical similarities of the lanthanides, chemical separation is not usually required for their applications, such as the addition of mischmetal to steel as an inclusion modifier to improve mechanical properties, or as catalysts for the cracking of petroleum. This property of cerium saved the life of writer Primo Levi at the Auschwitz concentration camp, when he found a supply of ferrocerium alloy and bartered it for food.
Pigments and phosphors
The photostability of pigments can be enhanced by the addition of cerium, as it provides pigments with lightfastness and prevents clear polymers from darkening in sunlight. An example of a cerium compound used on its own as an inorganic pigment is the vivid red cerium(III) sulfide (cerium sulfide red), which stays chemically inert up to very high temperatures. The pigment is a safer alternative to lightfast but toxic cadmium selenide-based pigments. The addition of cerium oxide to older cathode-ray tube television glass plates was beneficial, as it suppresses the darkening effect from the creation of F-center defects due to the continuous electron bombardment during operation. Cerium is also an essential component as a dopant for phosphors used in CRT TV screens, fluorescent lamps, and later white light-emitting diodes. The most commonly used example is cerium(III)-doped yttrium aluminium garnet (Ce:YAG) which emits green to yellow-green light (550–530 nm) and also behaves as a scintillator.
Other uses
Cerium salts, such as the sulfides Ce2S3 and Ce3S4, were considered during the Manhattan Project as advanced refractory materials for the construction of crucibles which could withstand the high temperatures and strongly reducing conditions when casting plutonium metal. Despite desirable properties, these sulfides were never widely adopted due to practical issues with their synthesis. Cerium is used as alloying element in aluminium to create castable eutectic aluminium alloys with 6–16 wt.% Ce, to which other elements such as Mg, Ni, Fe and Mn can be added. These Al-Ce alloys have excellent high temperature strength and are suitable for automotive applications (e.g. in cylinder heads). Other alloys of cerium include Pu-Ce and Pu-Ce-Co plutonium alloys, which have been used as nuclear fuel.
Other automotive applications for the lower sesquioxide are as a catalytic converter for the oxidation of CO and NO emissions in the exhaust gases from motor vehicles.
Biological role and precautions
The early lanthanides have been found to be essential to some methanotrophic bacteria living in volcanic mudpots, such as Methylacidiphilum fumariolicum: lanthanum, cerium, praseodymium, and neodymium are about equally effective. Cerium is otherwise not known to have biological role in any other organisms, but is not very toxic either; it does not accumulate in the food chain to any appreciable extent. Because it often occurs together with calcium in phosphate minerals, and bones are primarily calcium phosphate, cerium can accumulate in bones in small amounts that are not considered dangerous.
Cerium nitrate is an effective topical antimicrobial treatment for third-degree burns, although large doses can lead to cerium poisoning and methemoglobinemia. The early lanthanides act as essential cofactors for the methanol dehydrogenase of the methanotrophic bacterium Methylacidiphilum fumariolicum SolV, for which lanthanum, cerium, praseodymium, and neodymium alone are about equally effective.
Like all rare-earth metals, cerium is of low to moderate toxicity. A strong reducing agent, it ignites spontaneously in air at 65 to 80 °C. Fumes from cerium fires are toxic. Cerium reacts with water to produce hydrogen gas, and thus cerium fires can only be effectively extinguished using class D dry powder extinguishing media. Workers exposed to cerium have experienced itching, sensitivity to heat, and skin lesions. Cerium is not toxic when eaten, but animals injected with large doses of cerium have died due to cardiovascular collapse. Cerium is more dangerous to aquatic organisms because it damages cell membranes; it is not very soluble in water and can cause environmental contamination.
Cerium oxide, the most prevalent cerium compound in industrial applications, is not regulated in the United States by the Occupational Safety and Health Administration (OSHA) as a hazardous substance. In Russia, its occupational exposure limit is 5 mg/m. Elemental cerium has no established occupational or permissible exposure limits by the OSHA or American Conference of Governmental Industrial Hygienists, though it is classified as a flammable solid and regulated as such under the Globally Harmonized System of Classification and Labelling of Chemicals. Toxicological reports on cerium compounds have noted their cytotoxicity and contributions to pulmonary interstitial fibrosis in workers.
| Physical sciences | Chemical elements_2 | null |
953257 | https://en.wikipedia.org/wiki/Diving%20bell%20spider | Diving bell spider | The diving bell spider or water spider (Argyroneta aquatica) is the only species of spider known to live almost entirely under water. It is the only member of the genus Argyroneta. When out of the water, the spider ranges in colour from mid to dark brown, although the hairs on the abdomen give it a dark grey, velvet-like appearance. It is native to freshwater habitats in Europe and Asia.
Uniqueness of aquatic behavior
A. aquatica is the only known species of spider that spends almost all its life underwater, including resting, catching and eating prey, mating, egg laying, and overwintering. It only briefly surfaces to replenish its oxygen supply and occasionally will bring prey to the surface.
Several other spiders are semiaquatic, either periodically living underwater or willing to dive. For example, certain Desis species spend the high tide in an air-filled underwater retreat made from silk and forage on land in the intertidal zone during low tide. Some spiders living in periodically flooded habitats can survive for an extended period underwater by entering a coma-like state, up to 16–36 hours in Arctosa fulvolineata. Numerous species, including some Ancylometes, Dolomedes, Megadolomedes, Pardosa, Pirata, Thalassius and others, live above water at the surface, but may actively submerge for a prolonged period, are strong swimmers and will catch underwater prey. Several of these, as well as a few others, may dive into the water to avoid larger predators.
Distribution and habitat
A. aquatica is found in clean freshwater habitats with aquatic vegetation, such as lakes, ponds, canals, marshes and slow-moving streams. It ranges through much of mainland Europe (no records from Portugal, Greece and Albania), the British Isles and central to northern Asia ranging as far south as Iran and as far north as Siberia, up to latitude 62°N. Most of the range is inhabited by the nominate subspecies, but Japan has its own subspecies, the very similar A. a. japonica.
South Korea has a protected area, Yeoncheon Eundaeri water spider habitat, for this endangered species in that country.
Ecology
As with other spiders it breathes air; when submerged in water, an air bubble is trapped by a dense layer of hydrophobic hairs on its abdomen and legs, giving the abdomen a silvery appearance. The spider lives for about two years in captivity.
A. aquatica is able to remain submerged for prolonged periods of time due to the silk-based structure it constructs in order to retain an oxygen supply, named after the diving bell structure it resembles. The species range in size, although the size of females may be limited as they put more energy into building and maintaining their larger bells. Males are more active and on average almost 30% larger than females, measuring in head-and-body length compared to . This size differential favoring males is unusual for spiders, where sexual dimorphism is usually in favour of larger females. Theories suggest that the male's more active hunting style requires greater strength to overcome water resistance and counteract the buoyancy of their mobile air supplies. This larger body size is also associated with longer front legs, shown to affect diving ability and giving the males superiority in diving over the more sessile females. In addition, females build larger air bells than males, as the bell is also used to care for the offspring, forcing them to collect air from the surface much more frequently. And the larger the body, the more air is trapped on their abdomen, which means more buoyancy to overcome. Larger females therefore have higher energy costs than males of the same size and smaller females, which could limit the number of offspring they are able to produce.
The spiders prey on aquatic insects and crustaceans such as mosquito larvae and Daphnia. The spiders themselves fall prey to frogs and fish.
Diving bell
The appearance of the diving bell gave rise to the genus name Argyroneta, from the Greek "argyros" (ἄργυρος), meaning "silver", and "neta", a neologism (perhaps for *νητής) derived from the verb "neo" (νέω) "spin", intended to mean "spinner of silver". Both sexes build diving bell webs which are used for digesting prey, although only the female's larger bell is used for mating and raising offspring. Females spend most of their time within their bells, darting out to catch prey animals that touch the bell or the silk threads that anchor it and occasionally surfacing to replenish the air within the web. The bells built by males are typically smaller than females' and are replenished less often. It is thought that prior to mating, the male constructs a diving bell adjacent to the female's then spins a tunnel from his bell, breaking into hers to gain entrance. Mating takes place in the female's bell. The female spider then constructs an egg sac within her bell, laying between 30 and 70 eggs. Where this species moults is less clear, with some sources stating that it occurs below water in the diving bell and others that it occurs out of water.
Diving bells are irregularly constructed sheets of silk and an unknown protein-based hydrogel which is spun between submerged water plants then inflated with air brought down from the surface by the builder. Studies have considered gas diffusion between the diving bell and the spiders' aquatic environment. The silk is waterproof but allows gas exchange with the surrounding water. There is net diffusion of oxygen into the bell and net diffusion of carbon dioxide out. This process is driven by differences in partial pressure. The production of carbon dioxide and use of oxygen by the spider maintains the concentration gradient, required for diffusion. However, there is net diffusion of nitrogen out of the bell, resulting in a gradually shrinking air bubble which must be regularly replenished by the spider.
Larger spiders are able to produce larger bubbles which have a consequently higher oxygen conductance, but all spiders of this species are able to enlarge their bells in response to increased oxygen demands in low aquatic P(O2) environments. These spiders voluntarily tolerate internal conditions of low oxygen, enlarging their bells with air when the P(O2) drops below 1 kPa; this replenishment process may not need to occur for several days, in some cases. This system has been referred to as "the water spider's aqua-lung of air bubbles", though an aqua-lung lacks gas exchange with the surroundings; this system is more properly regarded as an inorganic form of gill.
Bite
Their bite is often described as being very painful to humans and as causing localised inflammation, vomiting, and slight feverishness that disappears after 5–10 days. However, solid evidence is lacking, with information being based on old and unverified reports because recent confirmed and published reports are lacking, leading some sources to refer to its bite as reputedly painful.
| Biology and health sciences | Spiders | Animals |
953369 | https://en.wikipedia.org/wiki/Virtual%20image | Virtual image | In optics, the image of an object is defined as the collection of focus points of light rays coming from the object. A real image is the collection of focus points made by converging rays, while a virtual image is the collection of focus points made by backward extensions of diverging rays. In other words, a virtual image is found by tracing real rays that emerge from an optical device (lens, mirror, or some combination) backward to perceived or apparent origins of ray divergences.
There is a concept virtual object that is similarly defined; an object is virtual when forward extensions of rays converge toward it. This is observed in ray tracing for a multi-lenses system or a diverging lens. For the diverging lens, forward extension of converging rays toward the lens will meet the converging point, so the point is a virtual object.
For a (refracting) lens, the real image of an object is formed on the opposite side of the lens while the virtual image is formed on the same side as the object. For a (reflecting) mirror, the real image is on the same side as the object while the virtual image is on the opposite side of, or "behind", the mirror. In diagrams of optical systems, virtual rays (forming virtual images) are conventionally represented by dotted lines, to contrast with the solid lines of real rays.
Because the rays never really converge, a virtual image cannot be projected onto a screen by putting it at the location of the virtual image. In contrast, a real image can be projected on the screen as it is formed by rays that converge on a real location. A real image can be projected onto a diffusely reflecting screen so people can see the image (the image on the screen plays as an object to be imaged by human eyes).
A plane mirror forms a virtual image positioned behind the mirror. Although the rays of light seem to come from behind the mirror, light from the source only exists in front of the mirror. The image in a plane mirror is not magnified (that is, the image is the same size as the object) and appears to be as far behind the mirror as the object is in front of the mirror.
A diverging lens (one that is thicker at the edges than the middle) or a concave mirror forms a virtual image. Such an image is reduced in size when compared to the original object. A converging lens (one that is thicker in the middle than at the edges) or a convex mirror is also capable of producing a virtual image if the object is within the focal length. Such an image will be magnified. In contrast, an object placed in front of a converging lens or concave mirror at a position beyond the focal length produces a real image. Such an image will be magnified if the position of the object is within twice the focal length, or else the image will be reduced if the object is further than this distance.
| Physical sciences | Optics | Physics |
953913 | https://en.wikipedia.org/wiki/Emperor%20scorpion | Emperor scorpion | The emperor scorpion (Pandinus imperator) is a species of scorpion native to rainforests and savannas in West Africa. It is one of the largest scorpions in the world and lives for six to eight years. Its body is black, but like other scorpions it glows pastel green or blue under ultraviolet light. It is a popular species in the pet trade, and is protected by CITES to prevent over-collecting that might affect the species' survival.
Description
The emperor scorpion (Pandinus imperator) is one of the largest species of scorpion in the world, with adults averaging about in length and a weight of 30 g. However, some species of forest scorpions are fairly similar to the emperor scorpion in size, and one scorpion, Heterometrus swammerdami, holds the record for being the world's largest scorpion at 9 inches (23 cm) in length. The large pincers are blackish-red and have a granular texture. The front part of the body, or prosoma, is made up of four sections, each with a pair of legs. Behind the fourth pair of legs are comb-like structures known as pectines, which tend to be longer in males than in females. The tail, known as the metasoma, is long and curves back over the body. It ends in the large receptacle containing the venom glands and is tipped with a sharp, curved stinger. Scorpion stings can be categorized as mild (similar to a bee sting) to severe to humans depending on the species. Most people are not severely affected by the emperor scorpion's sting, though some people may be allergic to scorpion stings in general. Sensory hairs cover the pincers and tail, enabling the emperor scorpion to detect prey through vibrations in the air and ground.
When gravid (pregnant), the body of a female expands to expose the whitish membranes connecting the segments. The emperor scorpion fluoresces greenish-blue under ultra-violet light.
They are known for their docile behavior and almost harmless sting; they do not use their sting to defend themselves when they are adults, however, they may use it in their adolescent stages. They prefer to use their pincers to crush and dismember their prey. Their exoskeleton is very sclerotic, causing them to have a metallic greenish-black color.
Emperor scorpions are often confused with a similar genus (Heterometrus), and are one of the most famous scorpions.
Different ion channel toxins have been isolated from the venom of the emperor scorpion, including Pi1, Pi2, Pi3, Pi4 and Pi7.
Habitat and distribution
The emperor scorpion is an African rainforest species, but also present in savanna. It is found in a number of African countries, including Benin, Burkina Faso, Côte d’Ivoire, Gambia, Ghana, Guinea, Guinea-Bissau, Togo, Liberia, Mali, Nigeria, Senegal, Sierra Leone and Cameroon.
This species inhabits both tropical forest and open savannas. The emperor scorpion burrows beneath the soil and hides beneath rocks and debris, and also often burrows in termite mounds.
Feeding habits
In the wild, emperor scorpions primarily consume insects and other terrestrial invertebrates, although termites constitute a large portion of their diet. Larger vertebrates, such as rodents and lizards, are occasionally eaten. Emperor scorpions will burrow through termite mounds up to 6 feet deep in order to hunt prey. Their large claws help in tearing apart prey while their tail stinger injects venom at the same time for liquifying food. Juveniles rely on their venomous sting to paralyze prey while adults use their large claws to tear apart prey.
Conservation and human impact
African emperor scorpion venom contains the toxins imperatoxin and pandinotoxin.
Due to its docile nature, large size, and hardiness, P. imperator is a popular scorpion in the pet trade, which has led to such over-collecting in the wild that it is now a CITES-listed animal. They feed readily on crickets and worms available to keepers, and they can live up to 8 years in captivity.
Despite their imposing size, emperor scorpion stings are usually mild, their venom does not cause severe symptoms in most people.
| Biology and health sciences | Scorpions | Animals |
954281 | https://en.wikipedia.org/wiki/Cache%20replacement%20policies | Cache replacement policies | In computing, cache replacement policies (also known as cache replacement algorithms or cache algorithms) are optimizing instructions or algorithms which a computer program or hardware-maintained structure can utilize to manage a cache of information. Caching improves performance by keeping recent or often-used data items in memory locations which are faster, or computationally cheaper to access, than normal memory stores. When the cache is full, the algorithm must choose which items to discard to make room for new data.
Overview
The average memory reference time is
where
= miss ratio = 1 - (hit ratio)
= time to make main-memory access when there is a miss (or, with a multi-level cache, average memory reference time for the next-lower cache)
= latency: time to reference the cache (should be the same for hits and misses)
= secondary effects, such as queuing effects in multiprocessor systems
A cache has two primary figures of merit: latency and hit ratio. A number of secondary factors also affect cache performance.
The hit ratio of a cache describes how often a searched-for item is found. More efficient replacement policies track more usage information to improve the hit rate for a given cache size.
The latency of a cache describes how long after requesting a desired item the cache can return that item when there is a hit. Faster replacement strategies typically track of less usage information—or, with a direct-mapped cache, no information—to reduce the time required to update the information. Each replacement strategy is a compromise between hit rate and latency.
Hit-rate measurements are typically performed on benchmark applications, and the hit ratio varies by application. Video and audio streaming applications often have a hit ratio near zero, because each bit of data in the stream is read once (a compulsory miss), used, and then never read or written again. Many cache algorithms (particularly LRU) allow streaming data to fill the cache, pushing out information which will soon be used again (cache pollution). Other factors may be size, length of time to obtain, and expiration. Depending on cache size, no further caching algorithm to discard items may be needed. Algorithms also maintain cache coherence when several caches are used for the same data, such as multiple database servers updating a shared data file.
Policies
Bélády's algorithm
The most efficient caching algorithm would be to discard information which would not be needed for the longest time; this is known as Bélády's optimal algorithm, optimal replacement policy, or the clairvoyant algorithm. Since it is generally impossible to predict how far in the future information will be needed, this is unfeasible in practice. The practical minimum can be calculated after experimentation, and the effectiveness of a chosen cache algorithm can be compared.
When a page fault occurs, a set of pages is in memory. In the example, the sequence of 5, 0, 1 is accessed by Frame 1, Frame 2, and Frame 3 respectively. When 2 is accessed, it replaces value 5 (which is in frame 1, predicting that value 5 will not be accessed in the near future. Because a general-purpose operating system cannot predict when 5 will be accessed, Bélády's algorithm cannot be implemented there.
Random replacement (RR)
Random replacement selects an item and discards it to make space when necessary. This algorithm does not require keeping any access history. It has been used in ARM processors due to its simplicity, and it allows efficient stochastic simulation.
Simple queue-based policies
First in first out (FIFO)
With this algorithm, the cache behaves like a FIFO queue; it evicts blocks in the order in which they were added, regardless of how often or how many times they were accessed before.
Last in first out (LIFO) or First in last out (FILO)
The cache behaves like a stack, and unlike a FIFO queue. The cache evicts the block added most recently first, regardless of how often or how many times it was accessed before.
SIEVE
SIEVE is a simple eviction algorithm designed specifically for web caches, such as key-value caches and Content Delivery Networks.
It uses the idea of lazy promotion and quick demotion. Therefore, SIEVE does not update the global data structure at cache hits and delays the update till eviction time; meanwhile, it quickly evicts newly inserted objects because cache workloads tend to show high one-hit-wonder ratios, and most of the new objects are not worthwhile to be kept in the cache. SIEVE uses a single FIFO queue and uses a moving hand to select objects to evict. Objects in the cache have one bit of metadata indicating whether the object has been requested after being admitted into the cache. The eviction hand points to the tail of the queue at the beginning and moves toward the head over time. Compared with the CLOCK eviction algorithm, retained objects in SIEVE stay in the old position. Therefore, new objects are always at the head, and the old objects are always at the tail. As the hand moves toward the head, new objects are quickly evicted (quick demotion), which is the key to the high efficiency in the SIEVE eviction algorithm. SIEVE is simpler than LRU, but achieves lower miss ratios than LRU on par with state-of-the-art eviction algorithms. Moreover, on stationary skewed workloads, SIEVE is better than existing known algorithms including LFU.
Simple recency-based policies
Least Recently Used (LRU)
Discards least recently used items first. This algorithm requires keeping track of what was used and when, which is cumbersome. It requires "age bits" for cache lines, and tracks the least recently used cache line based on these age bits. When a cache line is used, the age of the other cache lines changes. LRU is a family of caching algorithms, that includes 2Q by Theodore Johnson and Dennis Shasha and LRU/K by Pat O'Neil, Betty O'Neil and Gerhard Weikum. The access sequence for the example is A B C D E D F:
When A B C D is installed in the blocks with sequence numbers (increment 1 for each new access) and E is accessed, it is a miss and must be installed in a block. With the LRU algorithm, E will replace A because A has the lowest rank (A(0)). In the next-to-last step, D is accessed and the sequence number is updated. F is then accessed, replacing Bwhich had the lowest rank, (B(1)).
Time-aware, least-recently used
Time-aware, least-recently-used (TLRU) is a variant of LRU designed for when the contents of a cache have a valid lifetime. The algorithm is suitable for network cache applications such as information-centric networking (ICN), content delivery networks (CDNs) and distributed networks in general. TLRU introduces a term: TTU (time to use), a timestamp of content (or a page) which stipulates the usability time for the content based on its locality and the content publisher. TTU provides more control to a local administrator in regulating network storage.
When content subject to TLRU arrives, a cache node calculates the local TTU based on the TTU assigned by the content publisher. The local TTU value is calculated with a locally-defined function. When the local TTU value is calculated, content replacement is performed on a subset of the total content of the cache node. TLRU ensures that less-popular and short-lived content is replaced with incoming content.
Most-recently-used (MRU)
Unlike LRU, MRU discards the most-recently-used items first. At the 11th VLDB conference, Chou and DeWitt said: "When a file is being repeatedly scanned in a [looping sequential] reference pattern, MRU is the best replacement algorithm." Researchers presenting at the 22nd VLDB conference noted that for random access patterns and repeated scans over large datasets (also known as cyclic access patterns), MRU cache algorithms have more hits than LRU due to their tendency to retain older data. MRU algorithms are most useful in situations where the older an item is, the more likely it is to be accessed. The access sequence for the example is A B C D E C D B:
A B C D are placed in the cache, since there is space available. At the fifth access (E), the block which held D is replaced with E since this block was used most recently. At the next access (to D), C is replaced since it was the block accessed just before D.
Segmented LRU (SLRU)
An SLRU cache is divided into two segments: probationary and protected. Lines in each segment are ordered from most- to least-recently-accessed. Data from misses is added to the cache at the most-recently-accessed end of the probationary segment. Hits are removed from where they reside and added to the most-recently-accessed end of the protected segment; lines in the protected segment have been accessed at least twice. The protected segment is finite; migration of a line from the probationary segment to the protected segment may force the migration of the LRU line in the protected segment to the most-recently-used end of the probationary segment, giving this line another chance to be accessed before being replaced. The size limit of the protected segment is an SLRU parameter which varies according to I/O workload patterns. When data must be discarded from the cache, lines are obtained from the LRU end of the probationary segment.
LRU approximations
LRU may be expensive in caches with higher associativity. Practical hardware usually employs an approximation to achieve similar performance at a lower hardware cost.
Pseudo-LRU (PLRU)
For CPU caches with large associativity (generally > four ways), the implementation cost of LRU becomes prohibitive. In many CPU caches, an algorithm that almost always discards one of the least recently used items is sufficient; many CPU designers choose a PLRU algorithm, which only needs one bit per cache item to work. PLRU typically has a slightly-worse miss ratio, slightly-better latency, uses slightly less power than LRU, and has a lower overhead than LRU.
Bits work as a binary tree of one-bit pointers which point to a less-recently-used sub-tree. Following the pointer chain to the leaf node identifies the replacement candidate. With an access, all pointers in the chain from the accessed way's leaf node to the root node are set to point to a sub-tree which does not contain the accessed path. The access sequence in the example is A B C D E:
When there is access to a value (such as A) and it is not in the cache, it is loaded from memory and placed in the block where the arrows are pointing in the example. After that block is placed, the arrows are flipped to point the opposite way. A, B, C and D are placed; E replaces A as the cache fills because that was where the arrows were pointing, and the arrows which led to A flip to point in the opposite direction (to B, the block which will be replaced on the next cache miss).
Clock-Pro
The LRU algorithm cannot be implemented in the critical path of computer systems, such as operating systems, due to its high overhead; Clock, an approximation of LRU, is commonly used instead. Clock-Pro is an approximation of LIRS for low-cost implementation in systems. Clock-Pro has the basic Clock framework, with three advantages. It has three "clock hands" (unlike Clock's single "hand"), and can approximately measure the reuse distance of data accesses. Like LIRS, it can quickly evict one-time-access or low-locality data items. Clock-Pro is as complex as Clock, and is easy to implement at low cost. The buffer-cache replacement implementation in the 2017 version of Linux combines LRU and Clock-Pro.
Simple frequency-based policies
Least frequently used (LFU)
The LFU algorithm counts how often an item is needed; those used less often are discarded first. This is similar to LRU, except that how many times a block was accessed is stored instead of how recently. While running an access sequence, the block which was used the fewest times will be removed from the cache.
Least frequent recently used (LFRU)
The least frequent recently used (LFRU) algorithm combines the benefits of LFU and LRU. LFRU is suitable for network cache applications such as ICN, CDNs, and distributed networks in general. In LFRU, the cache is divided into two partitions: privileged and unprivileged. The privileged partition is protected and, if content is popular, it is pushed into the privileged partition. In replacing the privileged partition, LFRU evicts content from the unprivileged partition; pushes content from the privileged to the unprivileged partition, and inserts new content into the privileged partition. LRU is used for the privileged partition and an approximated LFU (ALFU) algorithm for the unprivileged partition.
LFU with dynamic aging (LFUDA)
A variant, LFU with dynamic aging (LFUDA), uses dynamic aging to accommodate shifts in a set of popular objects; it adds a cache-age factor to the reference count when a new object is added to the cache or an existing object is re-referenced. LFUDA increments cache age when evicting blocks by setting it to the evicted object's key value, and the cache age is always less than or equal to the minimum key value in the cache. If an object was frequently accessed in the past and becomes unpopular, it will remain in the cache for a long time (preventing newly- or less-popular objects from replacing it). Dynamic aging reduces the number of such objects, making them eligible for replacement, and LFUDA reduces cache pollution caused by LFU when a cache is small.
S3-FIFO
This is a new eviction algorithm designed in 2023. Compared to existing algorithms, which mostly build on LRU (least-recently-used), S3-FIFO only uses three FIFO queues: a small queue occupying 10% of cache space, a main queue that uses 90% of the cache space, and a ghost queue that only stores object metadata. The small queue is used to filter out one-hit-wonders (objects that are only accessed once in a short time window); the main queue is used to store popular objects and uses reinsertion to keep them in the cache; and the ghost queue is used to catch potentially-popular objects that are evicted from the small queue. Objects are first inserted into the small queue (if they are not found in the ghost queue, otherwise inserted into the main queue); upon eviction from the small queue, if an object has been requested, it is reinserted into the main queue, otherwise, it is evicted and the metadata is tracked in the ghost queue.
S3-FIFO demonstrates that FIFO queues are sufficient to design efficient and scalable eviction algorithms. Compared to LRU and LRU-based algorithms, S3-FIFO can achieve 6x higher throughput. Besides, on web cache workloads, S3-FIFO achieves the lowest miss ratio among 11 state-of-the-art algorithms the authors compared with.
RRIP-style policies
RRIP-style policies are the basis for other cache replacement policies, including Hawkeye.
Re-Reference Interval Prediction (RRIP)
RRIP is a flexible policy, proposed by Intel, which attempts to provide good scan resistance while allowing older cache lines that have not been reused to be evicted. All cache lines have a prediction value, the RRPV (re-reference prediction value), that should correlate with when the line is expected to be reused. The RRPV is usually high on insertion; if a line is not reused soon, it will be evicted to prevent scans (large amounts of data used only once) from filling the cache. When a cache line is reused the RRPV is set to zero, indicating that the line has been reused once and is likely to be reused again.
On a cache miss, the line with an RRPV equal to the maximum possible RRPV is evicted; with 3-bit values, a line with an RRPV of 23 - 1 = 7 is evicted. If no lines have this value, all RRPVs in the set are increased by 1 until one reaches it. A tie-breaker is needed, and usually, it is the first line on the left. The increase is needed to ensure that older lines are aged properly and will be evicted if they are not reused.
Static RRIP (SRRIP)
SRRIP inserts lines with an RRPV value of maxRRPV; a line which has just been inserted will be the most likely to be evicted on a cache miss.
Bimodal RRIP (BRRIP)
SRRIP performs well normally, but suffers when the working set is much larger than the cache size and causes cache thrashing. This is remedied by inserting lines with an RRPV value of maxRRPV most of the time, and inserting lines with an RRPV value of maxRRPV - 1 randomly with a low probability. This causes some lines to "stick" in the cache, and helps prevent thrashing. BRRIP degrades performance, however, on non-thrashing accesses. SRRIP performs best when the working set is smaller than the cache, and BRRIP performs best when the working set is larger than the cache.
Dynamic RRIP (DRRIP)
DRRIP uses set dueling to select whether to use SRRIP or BRRIP. It dedicates a few sets (typically 32) to use SRRIP and another few to use BRRIP, and uses a policy counter which monitors set performance to determine which policy will be used by the rest of the cache.
Policies approximating Bélády's algorithm
Bélády's algorithm is the optimal cache replacement policy, but it requires knowledge of the future to evict lines that will be reused farthest in the future. A number of replacement policies have been proposed which attempt to predict future reuse distances from past access patterns, allowing them to approximate the optimal replacement policy. Some of the best-performing cache replacement policies attempt to imitate Bélády's algorithm.
Hawkeye
Hawkeye attempts to emulate Bélády's algorithm by using past accesses by a PC to predict whether the accesses it produces generate cache-friendly (used later) or cache-averse accesses (not used later). It samples a number of non-aligned cache sets, uses a history of length and emulates Bélády's algorithm on these accesses. This allows the policy to determine which lines should have been cached and which should not, predicting whether an instruction is cache-friendly or cache-averse. This data is then fed into an RRIP; accesses from cache-friendly instructions have a lower RRPV value (likely to be evicted later), and accesses from cache-averse instructions have a higher RRPV value (likely to be evicted sooner). The RRIP backend makes the eviction decisions. The sampled cache and OPT generator set the initial RRPV value of the inserted cache lines. Hawkeye won the CRC2 cache championship in 2017, and Harmony is an extension of Hawkeye which improves prefetching performance.
Mockingjay
Mockingjay tries to improve on Hawkeye in several ways. It drops the binary prediction, allowing it to make more fine-grained decisions about which cache lines to evict, and leaves the decision about which cache line to evict for when more information is available.
Mockingjay keeps a sampled cache of unique accesses, the PCs that produced them, and their timestamps. When a line in the sampled cache is accessed again, the time difference will be sent to the reuse distance predictor. The RDP uses temporal difference learning, where the new RDP value will be increased or decreased by a small number to compensate for outliers; the number is calculated as . If the value has not been initialized, the observed reuse distance is inserted directly. If the sampled cache is full and a line needs to be discarded, the RDP is instructed that the PC that last accessed it produces streaming accesses.
On an access or insertion, the estimated time of reuse (ETR) for this line is updated to reflect the predicted reuse distance. On a cache miss, the line with the highest ETR value is evicted. Mockingjay has results which are close to the optimal Bélády's algorithm.
Machine-learning policies
A number of policies have attempted to use perceptrons, markov chains or other types of machine learning to predict which line to evict. Learning augmented algorithms also exist for cache replacement.
Other policies
Low inter-reference recency set (LIRS)
LIRS is a page replacement algorithm with better performance than LRU and other, newer replacement algorithms. Reuse distance is a metric for dynamically ranking accessed pages to make a replacement decision. LIRS addresses the limits of LRU by using recency to evaluate inter-reference recency (IRR) to make a replacement decision.
In the diagram, X indicates that a block is accessed at a particular time. If block A1 is accessed at time 1, its recency will be 0; this is the first-accessed block and the IRR will be 1, since it predicts that A1 will be accessed again in time 3. In time 2, since A4 is accessed, the recency will become 0 for A4 and 1 for A1; A4 is the most recently accessed object, and the IRR will become 4. At time 10, the LIRS algorithm will have two sets: an LIR set = {A1, A2} and an HIR set = {A3, A4, A5}. At time 10, if there is access to A4 a miss occurs; LIRS will evict A5 instead of A2 because of its greater recency.
Adaptive replacement cache
Adaptive replacement cache (ARC) constantly balances between LRU and LFU to improve the combined result. It improves SLRU by using information about recently-evicted cache items to adjust the size of the protected and probationary segments to make the best use of available cache space.
Clock with adaptive replacement
Clock with adaptive replacement (CAR) combines the advantages of ARC and Clock. CAR performs comparably to ARC, and outperforms LRU and Clock. Like ARC, CAR is self-tuning and requires no user-specified parameters.
Multi-queue
The multi-queue replacement (MQ) algorithm was developed to improve the performance of a second-level buffer cache, such as a server buffer cache, and was introduced in a paper by Zhou, Philbin, and Li. The MQ cache contains an m number of LRU queues: Q0, Q1, ..., Qm-1. The value of m represents a hierarchy based on the lifetime of all blocks in that queue.
Pannier
Pannier is a container-based flash caching mechanism which identifies containers whose blocks have variable access patterns. Pannier has a priority-queue-based survival-queue structure to rank containers based on their survival time, which is proportional to live data in the container.
Static analysis
Static analysis determines which accesses are cache hits or misses to indicate the worst-case execution time of a program. An approach to analyzing properties of LRU caches is to give each block in the cache an "age" (0 for the most recently used) and compute intervals for possible ages. This analysis can be refined to distinguish cases where the same program point is accessible by paths that result in misses or hits. An efficient analysis may be obtained by abstracting sets of cache states by antichains which are represented by compact binary decision diagrams.
LRU static analysis does not extend to pseudo-LRU policies. According to computational complexity theory, static-analysis problems posed by pseudo-LRU and FIFO are in higher complexity classes than those for LRU.
| Mathematics | Algorithms | null |
956001 | https://en.wikipedia.org/wiki/Gonepteryx%20rhamni | Gonepteryx rhamni | Gonepteryx rhamni, commonly named the common brimstone, is a butterfly of the family Pieridae. It lives throughout the Palearctic zone and is commonly found across Europe, Asia, and North Africa. Across much of its range, it is the only species of its genus, and is therefore simply known locally as the brimstone. Its wing span size is . It should not be confused with the brimstone moth Opisthograptis luteolata.
The brimstone relies on two species of buckthorn plants as host plants for its larvae; this influences its geographic range and distribution, as these plants are commonly found in wetlands. The adult brimstone travels to woodland areas to spend seven months overwintering. In spring when their host plants have developed, they return to the wetlands to breed and lay eggs. Both the larval and adult forms of the common brimstone have protective coloration and behaviour that decreases their chances of being recognised and subsequently preyed upon.
The adult common brimstone has sexual dimorphism in its wing coloration: males have yellow wings and iridescence while females have greenish-white wings and are not iridescent. This iridescence is affected by environmental factors.
Taxonomy
It was first described and published in Linnaeus's book, the 10th edition of Systema Naturae in 1758.
Brimstone is an old name for sulphur, the colour which matches the colour of the male's wings.
Distribution and habitat
The common brimstone can be commonly found throughout the Palearctic. Individuals have been seen from western Europe to east Asia. The high mobility of this butterfly allows it to search widely for new host plant locations and expand its range. While the geographic distribution of the adult is larger than that of its host plant, its range is nevertheless limited by the presence of host plants due to the needs of its larval stage.
The common brimstone uses various environments for different stages of its life cycle. The butterfly inhabits wetlands during mating and breeding season, as they provide ideal areas for oviposition due to an abundance of host plants like the alder buckthorn. The common brimstone prefers laying eggs on younger host plants with late bud-bursts that are isolated from other plants in the area and exposed to both open space and sun. During the winter, adult brimstones travel to woodlands to hibernate, as they provide ideal overwintering sites with shelters such as evergreen foliage and holly. The common brimstone has an appearance that is highly similar to the leaves of these plants, so during hibernation it can remain hidden. In other seasons, habitat selection is also affected by the abundance of nectar as a food source for adult brimstones.
Food resources
Caterpillar
Larval brimstones appear to feed on only two plant sources: the alder buckthorn (Frangula alnus) and the common buckthorn (Rhamnus carthartica). This influences the distribution of the adult brimstone, as the presence of these two buckthorn species is necessary for the survival of their offspring.
Adult
Unlike their larval forms, which are specialised for particular host plants, adult brimstones are not specialised nectar feeders. The common brimstone heavily feeds on the nectar of several flowering species including knapweed (Centaurea jacea) and scabious (Knautia arvensis and Succisa pratensis). However, brimstones have also been observed feeding on the nectar of coltsfoot (Tussilago farfara) in April and May and have been recorded gathering nectar from many other species of flowers. Adult food plant availability is another factor that is important for habitat selection.
Parental care
Oviposition
The common brimstone is univoltine, meaning that it lays one generation of eggs each year. There are several ideal characteristics of the particular host plants chosen for oviposition. Adult brimstones lay eggs on the underside of the leaves of the two species of host plants, where they are less conspicuous. The high mobility of G. rhamni enables the butterflies to find even the most isolated host plants in an area, which are more ideal for their offspring. Eggs are more likely to be deposited on outlying plants, leading to reduced vulnerability as fewer predators are attracted to these plants. Another factor is damage; undamaged plants indicate the absence of other eggs, as brimstone larvae leave holes in the leaves of the plants on which they feed. Since predators and parasites are attracted to damaged plants through chemical or visual signals, less damage leads to greater offspring survival since eggs are less likely to be detected. Plants exposed to both sunlight and the open lead to reduced chances of predation and parasitism as well, and are more accessible to adult butterflies. Larvae can also benefit from decreased host plant defences; juvenile plants and plants with late bud-bursts produce fewer toxic defence chemicals, as resources are directed more towards plant growth.
Life cycle
The common brimstone is one of the longest-living butterflies, with a life expectancy ranging from 10 months to a year. Due to its hibernation and life cycle, it has one generation per year. Development from the laid egg to the emergence of the imago is approximately 50 days. However, the adult brimstone spends a large portion of its life in an overwintering state. The brimstone is highly mobile, feeding and travelling to regions ideal for hibernation during the late summer and fall, and returning to regions ideal for mating and egg-laying during the spring.
Egg
Adult common brimstones lay eggs singly on the underside of buckthorn leaves. The eggs are around 1.3 mm tall, and are spindle-shaped in appearance. The eggs change colour over time, initially having a greenish-white colouration, then progressively darker shades of yellow, and finally brown before hatching.
Caterpillar
The larvae of the common brimstone undergo five instars, initially having a length of 1.7 mm in the first instar and reaching up to 34.9 mm in length when fully grown. The caterpillars have a green colouration with white hairs and dark tubercules across its length. When they first hatch, they move to the top side of the leaves and eat them, leaving characteristic hole patterns in their host plants. During the day, they feed and then rest in the open, lying still on the midrib of leaves, where their colouration makes them difficult to distinguish.
Pupa
Pupation occurs over approximately two weeks. The pupae are 22.2-23.8 mm in length and have the appearance of a curled leaf, with pointed ends and bulges in the middle. The pupae are secured to stems and leaves using silk; a cremastral hook attaches to a silk padding, and a length of silk secures the pupae around its middle. The pupae have a primarily green colouration, but right before adult emergence for males, the wing areas turn yellow.
Adult
Adults emerge during the summer, from June to August, and continue to feed until September. The common brimstone hibernates for the next seven months of winter, remaining inactive until April, where they then emerge and proceed to reproduce and lay eggs. Adult brimstones are highly abundant for several months after their emergence from overwintering. The common brimstone has sexual dichromism, with males having a sulphur yellow wing colouration and females having a greenish-white wing colouration. Additionally, males have iridescent dorsal wings that change in colour and appearance under ultraviolet light, while females do not. Both males and females have orange spots in the discoidal cell of each wing, pink head and antennae, and a thorax covered in white hair.
Migration
The common brimstone undergoes some regional migration between hibernation and breeding areas throughout the year, as seen in the different chemical composition of butterflies across varying seasons and regions. In general, there is movement towards wetlands to reproduce. After the eggs hatch, develop, and pupate, newly hatched adult butterflies emerge and disperse locally into both woodlands and wetlands to overwinter. Butterflies travel to the woodlands for overwintering, and no mating appears to occur within these habitats. Overwintering also occurs in the wetlands, where the host plant alder buckthorn is abundant. After emerging from overwintering, adult brimstones that were previously in the wetlands are joined by those that hibernated in woodlands, and the population breeds and lays eggs.
The environmental conditions of a particular year also affect migration, as seen in the elevational migrations of the common brimstone. Uphill migration is potentially influenced by habitat limitations, such as a lack of the forest cover that is required during overwintering. Brimstones travel to higher elevations for greater forest cover and reduced exposure to higher temperatures during their activities. Downhill migration is influenced by the need for larval resources such as host plants during breeding seasons - the butterflies travel to lower elevations in search for regions containing these plants, with adults commonly returning to the areas where they had been bred due to their long lifespan.
Enemies
Predators
Like most woodland Lepidoptera, G. rhamni is preyed upon by many species of birds and wasps. Both larvae and adult brimstones fall victim to predation and use means such as protective coloration and mimicry to avoid this fate.
Parasites
The common brimstone has two recorded species of parasites: the braconids Cotesia gonopterygis and Cotesia risilis. These two species of parasitoid wasps are completely specialised for G. rhamni, possibly due to the wide distribution of the butterfly and the host plants in its habitats. The broad presence of its host allows the wasps to be host-specific. The wasps are primarily associated with the presence of the food plant Frangula alnus due to its association with their host.
Protective colouration and behaviour
Both the larval and adult common brimstone exhibit cryptic colouration, meaning they match the colour of their habitats. Larvae are so difficult to see due to this colouration that they can remain in the open undetected. When not eating, the caterpillars remain still in a position alongside the midrib of leaves, making them even more difficult to spot. Adult brimstones are leaf-mimics, as they share similarities in shape, colour, and pattern to leaves. This allows them to blend in with their surroundings during vulnerable times like diapause (hibernation). When picked up, the butterflies become rigid and hide their legs from view in order to decrease their chances of being recognised.
Genetics of colour patterns
Pigmentation and structural coloration
Variation in coloration of Lepidoptera wings is caused by different structural and pigment components. These differences cause light to scatter in different ways, leading to the different colours. In the common brimstone, wing scales scatter light incoherently due to ovoid-shaped structures called beads that contain pigments. Due to these pigments, the beads absorb short wavelength light and scatter longer wavelengths outside of the pigment absorption spectrum, such as light in the complementary wavelength range.
Through chemical extraction and analysis, two possible pigments have been identified that may contribute to the common brimstone's wing coloration. Xanthopterin is responsible for the sulphur yellow colour of the male wings since it absorbs in the violet range. Leucopterin was extracted from the white wings of females. The difference in wing pigmentation contributes to the distinct iridescence patterns of males and females. Iridescence occurs due to visualised changes in coloration from the scattering of ultraviolet light. A male-only pattern of coloration due to this iridescence is seen exclusively under ultraviolet light, since females absorb light on the ultraviolet spectra. The presence of exclusively leucopterin in female wings explains the lack of iridescence in female common brimstones, since leucopterin absorbs only in the ultraviolet range. Therefore, the wings do not reflect and consequently do not scatter any ultraviolet light like male wings do.
In males, iridescence is indicated in that the wing pattern appears to visually change depending on the position of the ultraviolet light shone onto the wing. At some angles, a male pattern is seen, while at other angles, a female lack of pattern is seen. This is referred to as the "gynandromorphic effect". This demonstrates that the pattern appears to be optical, rather than pigmental, as the effect is only seen at certain angles and distances of light and changes with positions. If it were pigmental, these changes would not cause differences in iridescence.
The structural coloration of the male dorsal wings is affected by environmental factors. There is an increase in ultraviolet coloration coverage with increasing temperature, increasing precipitation, and decreasing latitude. This has been possibly attributed to several factors, such as the greater abundance and quality of resources in areas with these environmental conditions. Other possibilities include a better ability to assimilate resources as an indication of male mate quality. Because ultraviolet coloration is energetically expensive to develop, it could signal high male quality.
Mating
After the common brimstone emerges from hibernation, it travels towards habitats that contain larval host plants and mates. The brimstone is primarily monandrous, as demonstrated by the presence of usually only a single spermatophore in females throughout the mating season. Pairs are formed after the butterflies have engaged in a dalliance flight for a period of time. When a pair settles to mate, they do not take flight during copulation and remain paired for a long time of up to forty-eight hours.
Physiology
Vision
The common brimstone appears to have an innate preference for certain colours in nectar plants – red and blue inflorescences are common in heavily used nectar sources in some regions. G. rhamni also has a stronger reliance on visual indications such as colour compared with other butterfly species, which rely more on odour.
Olfaction
The common brimstone has an antennal response to the floral scent compounds of nectar plants, where neural activity in antennal olfactory receptors occurs in the presence of certain compounds. Research suggests that there are antennal olfactory receptors for phenylacetaldehyde and the terpene compounds oxoisophoroneoxide, oxoisophorone, and dihydrooxoisophorone, as these compounds elicited some of the strongest electrophysiological responses whether they were presented in natural or synthetic mixes of floral compounds. Additionally, these two compounds are present in the largest quantities in the nectar plants utilised by the brimstone, indicating that scent detection could be important for detecting food sources. This would contribute to more efficient foraging in adult butterflies, as odour could act as a cue for finding and distinguishing nectar plants, allowing more energy to be utilised for other activities such as reproduction.
Diapause
The adult common brimstone overwinters for seven months, remaining hidden and motionless throughout its hibernation. While both sexes have similar egg to adult development times, they differ in the times that they reach sexual maturity. The reproductive development of males begins just after pupal emergence, and continues during hibernation, which indicates that males may not be able to reproduce until after overwintering. For females, eggs remain undeveloped as the butterflies overwinter, and no reproductive development occurs until after emergence from hibernation.
The sexes also differ in times of emergence after overwintering. Emergence is correlated with temperature and hours of sunlight; a certain amount of both is necessary for the butterfly to emerge from hibernation and therefore influences when diapause ends. Males emerge earlier than females, as they are more willing to fly in lower temperatures than females. Since the common brimstone most closely follows monandrous mating patterns, males may emerge earlier to increase the number of mating chances and therefore reproductive success, as older males have had more time to develop and therefore have a greater advantage. In contrast, females emerge late due to the late seasonal development of host plants such as the alder buckthorn, since these plants are necessary for egg-laying. Female emergence is correlated with host plant development.
Conservation
As of 2010, G. rhamni does not appear to have a threatened conservation status according to IUCN standards. However, the butterfly has experienced significant population and distribution reduction in areas such as the Netherlands, where its numbers have declined to the point that based on IUCN criterion, it has reached endangered species status. The causes of this population decline are not fully determined, but there are several possible factors. Since the common brimstone is univoltine, it may have difficulties adapting to changing environmental conditions compared to species that have multiple generations a year. For example, there has been a decrease in suitable overwintering environments for the butterflies, with open woodland decreasing in favour of more urban areas. Nitrogen pollution, declining nectar supplies, and rapid ecological changes have also been suggested as other hypothetical factors. Concerns have been raised about the possible future increase of this population decline, but the butterfly mostly does not appear to be a conservation concern due to its widespread and common geographic presence.
| Biology and health sciences | Lepidoptera | Animals |
956656 | https://en.wikipedia.org/wiki/Maar | Maar | A maar is a broad, low-relief volcanic crater caused by a phreatomagmatic eruption (an explosion which occurs when groundwater comes into contact with hot lava or magma). A maar characteristically fills with water to form a relatively shallow crater lake, which may also be called a maar.
Maars range in size from across and from deep. Most maars fill with water to form natural lakes. Most maars have low rims composed of a mixture of loose fragments of volcanic rocks and rocks torn from the walls of the diatreme.
Etymology
The name maar comes from a Moselle Franconian dialect word used for the circular lakes of the Daun area of Germany. The word evolved from its first use in German in the modern geological sense in 1819 and is now used in English and in the geological sciences as the term for the explosion crater, even if water from rainfall might always have drained from the crater after the formation event. This extension in meaning was due to recognising that the lake may no longer exist. Since maar lakes are formed after initially ground or subsurface water interacts with a magma intrusion to create an explosion crater, the name came to be used for the crater type as well. The present definition of the term relates to both its common and scientific discourse use in language over two centuries. Depending upon context there may be other descriptors available to use in the geological sciences such as the term tuff ring or maar-diatreme volcanoes. These last are volcanoes produced by explosive eruptions that cut deeply into the country rock with the maar being "the crater cut into the ground and surrounded by an ejecta ring". A 2011 geological clarification of a maar is "Maar volcanoes are distinguished from other small volcanoes in having craters with their floor lying below the pre-eruptive surface".
Maar lakes and dry maars
Maar lakes, also referred to simply as maars, occur when groundwater or precipitation fills the funnel-shaped and usually round hollow of the maar depression formed by volcanic explosions. Examples of these types of maar are the three maars at Daun in the Eifel mountains of Germany.
A dry maar results when a maar lake dries out, becomes aggraded or silted up. An example of the latter is the Eckfelder Maar. Near Steffeln is the Eichholzmaar (also called the Gussweiher) which has dried out during the last century and is being renaturalised into a maar. In some cases the underlying rock is so porous that maar lakes are unable to form. After winters of heavy snow and rainfall many dry maars fill partially and temporarily with water; others contain small bogs or often artificial ponds that, however, only occupy part of the hollow.
Distribution
The largest known maars are found at Espenberg on the Seward Peninsula in northwest Alaska. These maars range in size from in diameter and a depth up to . These eruptions occurred in a period of about 100,000 years, with the youngest (the Devil Mountain Maar) occurring about 17,500 years ago. Their large size is due to the explosive reaction that occurs when magma comes into contact with permafrost. Hydromagmatic eruptions are increasingly explosive when the ratio of water to magma is low. Since permafrost melts slowly, it provides a steady source of water to the eruption while keeping the water to magma ratio low. This produces the prolonged, explosive eruptions that created these large maars. Examples of the Seward Peninsula maars include North Killeak Maar, South Killeak Maar, Devil Mountain Maar and Whitefish Maar.
Maars occur in western North America, Patagonia in South America, the Eifel region of Germany (where they were originally described), and in other geologically young volcanic regions of Earth. Elsewhere in Europe, La Vestide du Pal, a maar in the Ardèche department of France, is easily visible from the ground or air. Kilbourne Hole and Hunt's Hole, in southern New Mexico near El Paso, Texas, are maars. The Crocodile Lake in Los Baños in the Philippines, though originally thought to be a volcanic crater, is a maar. The carbon dioxide-saturated Lake Nyos in northwestern Cameroon is another example, as is Zuñi Salt Lake in New Mexico, a shallow saline lake that occupies a flat-floored crater about across and deep. Its low rim is composed of loose pieces of basaltic lava and wall rocks (sandstone, shale, limestone) of the underlying diatreme, as well as chunks of ancient crystalline rocks blasted upward from great depths. Maars in Canada are found in the Wells Gray-Clearwater volcanic field of east-central British Columbia and in kimberlite fields throughout Canada. Another field of maars is found in the Pali-Aike Volcanic Field in Patagonia, South America. and in the Sudanese Bayuda Volcanic Field. The Auckland volcanic field in the urban area of Auckland, New Zealand, has several maars, including the readily accessible Lake Pupuke in the North Shore suburb of Takapuna.
Arizona's Meteor Crater was for many years thought to be a maar of volcanic origin but it is now known to be an impact crater.
Examples
Germany
Eifel maars
In the Volcanic Eifel there are about 75 maars. Both lake-filled and dry maars are typical, though the latter are more common. The last eruptions took place at least 11,000 years ago, and many maars are older, as evidenced by their heavy erosion and less obvious shapes and volcanic features.
Water-filled maars of the Eifel
Dry maars of the Eifel
In the Eifel and Volcanic Eifel there are numerous dry maars:
Mosbrucher Weiher (4 km SE of Kelberg)
Booser Doppelmaar (W of Boos; near Kelberg)
Dreiser Weiher (W of Dreis-Brück, N of Daun)
Dürres Maar (SW of Gillenfeld)
Duppacher Weiher (near Duppach, NW of Gerolstein)
Geeser Maar (E of Gerolstein, N of Gees)
Eckfelder Maar (near Eckfeld)
Eigelbacher Maar (near Kopp, county of Daun; maar basin: c. 1200 m × 1200 m)
Hitsche Maar (NW of the Dürre Maar, smallest Eifel maar; Ø = 60 m)
Immerather Risch (Middle Low German: risch = reed bed; N of the Immerather Maar)
Gerolsteiner Maar (NE of Gerolstein)
Schalkenmehrener Maar E (of Schalkenmehren)
Schönfelder Maar (SW of Stadtkyll-Schönfeld)
Steffelner Laach or "Laach Maar" (near Steffeln)
Dehner Maar (near Reuth)
Walsdorfer Maar ("Schilierwiese"; S of Walsdorf; maar basin: c. 1150 m × 1000 m)
Wollmerather Maar (near Wollmerath)
Broader use of the term maar
The following volcanic features are often colloquially referred to as a "maar" or "maar lake", although they are not, strictly speaking, maars:
Windsborn Crater Lake and Hinkelsmaar in the Manderscheid Volcano Group near Bettenfeld, crater lakes of the Mosenberg (Eifel)|Mosenberg
Laacher See near Maria Laach, lake in a caldera of the Laacher See volcano
Strohner Märchen (south of the Pulvermaar), volcanic pipe which has become a maar
Papenkaule, a volcanic crater, and the associated eruption site of the Hagelskaule
Elfenmaar near Bad Bertrich, an almost entirely eroded stratovolcano
Rodder Maar near Niederdürenbach, the origin of which is unclear
Maars outside the Eifel
In Germany there are also several maars outside of the Eifel. A well-known example is the Messel pit, a former maar lake near Messel in the county of Darmstadt-Dieburg and which is known for its well preserved fossils. In addition in the Swabian Jura and the Albvorland (the Swabian Volcano) there are maar-forming volcanoes. Because the over 350 eruption points were only active in the Upper Miocene 17 to 11 million years ago, all the maars, apart from the dry maar, Randecker Maar and the Molach, are only detectable geologically. In the Ore Mountains near Hammerunterwiesenthal, the Hammerunterwiesenthal Maar formed about 30 million years ago during the Oligocene; the maar measures 2 kilometres from east to west and 1.4 kilometres from north to south.
Rest of Europe
The Chaîne des Puys in France contains numerous maars; Lake Albano in the Alban Mountains is a complex maar, and there is also a submarine maar (Kolumbo) near Santorini in Greece. The Campo de Calatrava Volcanic Field in Spain contains numerous maars; a typical example being the maar of Hoya del Mortero at Poblete in the Province of Ciudad Real.
Active maars were commonplace in Fife and Lothian, Scotland during the Carboniferous period. The location of one such maar was Elie Ness.
Americas
Active maar volcanoes are mainly known outside Europe.
In the US there are numerous maar areas, such as in Alaska (Ukinrek maars, Nunivak in the Bering Sea); in Washington (Battle Ground Lake); in Oregon (Fort Rock basin with the maars of Big Hole, Hole-in-the-Ground, Table Rock); in Death Valley National Park, California (Ubehebe Crater); in Nevada (Soda Lakes); as well as the maars of the White Rock Canyon, Mount Taylor, the Potrillo volcanic fields (Kilbourne Hole and Hunt's Hole), and Zuñi Salt Lake in New Mexico.
In Central Mexico, the Tarascan volcanic field contains several maars in the states of Michoacán and Guanajuato. In Nicaragua is the maar of Laguna de Xiloa, part of the Apoyeque volcano. From South America, there are known maars in Chile (e.g. Cerro Overo and Cerro Tujle in northern Chile). Jayu Khota is a maar in Bolivia.
Middle East and Africa
The maar of Lake Ram lies on the Golan Heights; further south maars occur in Africa (Bilate Volcanic Field and Haro Maja in the Butajiri-Silti-Volcanic Field, Ethiopia, the Bayuda Volcanic Field in the Sudan and Lake Nyos in the Oku Volcanic Field in Cameroon). In Saudi Arabia the Al Wahbah crater formed as a result of a maar eruption.
Asia and Oceania
In Japan there are maars in the Kirishima-Yaku volcanic field in the Kirishima-Yaku National Park on Kyushu. These include the several maars of the Ibusuki volcanic field such as Lake Unagi. On Honshu in Myōkō-Togakushi Renzan National Park there is Kagamiike Pond as well as many on the volcanic island of Miyake-jima, Izu Islands (Furumio, Mi'ike, Mizutamari, Shinmio).
Koranga Maar and Numundo Maar are in Papua New Guinea. Kawah Masemo maar is on Mount Sempu volcano in Indonesia. The San Pablo Volcanic Field in the Province of Laguna on the island of Luzon in the Philippines contains maars.
The Newer Volcanics Province in the States of South Australia and Victoria, Australia, has numerous maars, such as Mount Gambier, Mount Schank and Tower Hill, whose complex system of nested maars is enclosed by one of the largest maars in the world.
Foulden Maar in Otago, New Zealand, is an important fossil site, but there are many more maars in New Zealand. As already mentioned these include Lake Pupuke, but the Auckland volcanic field has other easily accessible maars such as the Mangere Lagoon, Orakei Basin, Panmure Basin, and Pukaki Lagoon. Elsewhere a recent example, only 4000 years old, is Lake Rotokawau in the Bay of Plenty Region.
Gallery
| Physical sciences | Volcanic landforms | Earth science |
956704 | https://en.wikipedia.org/wiki/Stopwatch | Stopwatch | A stopwatch is a timepiece designed to measure the amount of time that elapses between its activation and deactivation.
A large digital version of a stopwatch designed for viewing at a distance, as in a sports stadium, is called a stop clock. In manual timing, the clock is started and stopped by a person pressing a button. In fully automatic time, both starting and stopping are triggered automatically, by sensors.
The timing functions are traditionally controlled by two buttons on the case. Pressing the top button starts the timer running, and pressing the button a second time stops it, leaving the elapsed time displayed. A press of the second button then resets the stopwatch to zero. The second button is also used to record split times or lap times. When the split time button is pressed while the watch is running it allows the elapsed time to that point to be read, but the watch mechanism continues running to record total elapsed time. Pressing the split button a second time allows the watch to resume display of total time.
Mechanical stopwatches are powered by a mainspring, which must be wound up by turning the knurled knob at the top of the stopwatch.
Digital electronic stopwatches are available which, due to their crystal oscillator timing element, are much more accurate than mechanical timepieces. Because they contain a microchip, they often include date and time-of-day functions as well. Some may have a connector for external sensors, allowing the stopwatch to be triggered by external events, thus measuring elapsed time far more accurately than is possible by pressing the buttons with one's finger. The first digital timer used in organized sports was the Digitimer, developed by Cox Electronic Systems, Inc. of Salt Lake City Utah (1962). It utilized a Nixie-tube readout and provided a resolution of 1/1000 second. Its first use was in ski racing but was later used by the World University Games in Moscow, Russia, the U.S. NCAA, and in the Olympic trials.
The device is used when time periods must be measured precisely and with a minimum of complications. Laboratory experiments and sporting events like sprints are good examples.
The stopwatch function is also present as an additional function of many electronic devices such as wristwatches, cell phones, portable music players, and computers.
Humans are prone to make mistakes every time they use one. Normally, humans will take about 180–200 milliseconds to detect and respond to visual stimulus. However, in most situations where a stopwatch is used, there are indicators that the timing event is about to happen, and the manual action of starting/stopping the timer can be much more accurate. The average measurement error using manual timing was evaluated to be around 0.04 s when compared to electronic timing, in this case for a running sprint.
To get more accurate results, most researchers use the propagation of uncertainty equation in order to reduce any error in experiments.
is the sum of the uncertainty between and
is the value which is actually found from the experiment.
is the value of the uncertainty.
For example: If the result from measuring the width of a window is 1.50 ± 0.05 m, 1.50 will be and 0.05 will be .
Unit
In most science experiments, researchers will normally use SI or the International System of Units on any of their experiments. For stopwatches, the units of time that are generally used when observing a stopwatch are minutes, seconds, and 'one-hundredth of a second'.
Many mechanical stopwatches are of the 'decimal minute' type. These split one minute into 100 units of 0.6s each. This makes addition and subtraction of times easier than using regular seconds.
Types of stopwatches
| Technology | Clocks | null |
957081 | https://en.wikipedia.org/wiki/Oseltamivir | Oseltamivir | Oseltamivir, sold under the brand name Tamiflu among others, is an antiviral medication used to treat and prevent influenza A and influenza B, viruses that cause the flu. Many medical organizations recommend it in people who have complications or are at high risk of complications within 48 hours of first symptoms of infection. They recommend it to prevent infection in those at high risk, but not the general population. The Centers for Disease Control and Prevention (CDC) recommends that clinicians use their discretion to treat those at lower risk who present within 48 hours of first symptoms of infection. It is taken by mouth, either as a pill or liquid.
Recommendations regarding oseltamivir are controversial as are criticisms of the recommendations. A 2014 Cochrane Review concluded that oseltamivir does not reduce hospitalizations, and that there is no evidence of reduction in complications of influenza. Two meta-analyses have concluded that benefits in those who are otherwise healthy do not outweigh its risks. They also found little evidence regarding whether treatment changes the risk of hospitalization or death in high risk populations. However, another meta-analysis found that oseltamivir was effective for prevention of influenza at the individual and household levels.
Common side effects include vomiting, diarrhea, headache, and trouble sleeping. Other side effects may include psychiatric symptoms and seizures. In the United States it is recommended for influenza infection during pregnancy. It has been taken by a small number of pregnant women without signs of problems. Dose adjustment may be needed in those with kidney problems.
Oseltamivir was approved for medical use in the US in 1999. It was the first neuraminidase inhibitor available by mouth. It is on the World Health Organization's List of Essential Medicines but was downgraded to "complementary" status in 2017. A generic version was approved in the US in 2016. In 2022, it was the 205th most commonly prescribed medication in the United States, with more than 1million prescriptions.
Medical use
Oseltamivir is used for the prevention and treatment of influenza caused by influenza A and B viruses. It is on the World Health Organization's List of Essential Medicines. The WHO supports its use for severe illness due to confirmed or suspected influenza virus infection in critically ill people who have been hospitalized. Oseltamivir's risk-benefit ratio is controversial. In 2017, it was moved from the core to the complementary list based on its lower cost-effectiveness. The Expert Committee did not recommend the deletion of oseltamivir from the EML and EMLc, recognizing that it is the only medicine included on the Model Lists for critically ill patients with influenza and for influenza pandemic preparedness. However, the Committee noted that, since the inclusion of oseltamivir on the Model List in 2009, new evidence in seasonal and pandemic influenza has lowered earlier estimates of the magnitude of effect of oseltamivir on relevant clinical outcomes. The Committee recommended that the listing of oseltamivir be amended, moving the medicine from the core to the Complementary List, and that its use be restricted to severe illness due to confirmed or suspected influenza virus infection in critically ill hospitalized patients. The Expert Committee noted that WHO guidelines for pharmacological management of pandemic and seasonal influenza would be updated in 2017: unless new information is provided to support the use of oseltamivir in seasonal and pandemic outbreaks, the next Expert Committee might consider oseltamivir for deletion.
High-risk people
The US Centers for Disease Control and Prevention (CDC), European Centre for Disease Prevention and Control (ECDC), Public Health England and the American Academy of Pediatrics (AAP) recommend the use of oseltamivir for people who have complications or are at high risk for complications. This includes those who are hospitalized, young children, those over the age of 65, people with other significant health problems, those who are pregnant, and Indigenous peoples of the Americas among others. The Infectious Disease Society of America takes the same position as the CDC.
A systematic review of systematic reviews in PLoS One did not find evidence for benefits in people who are at risk, noting that "the trials were not designed or powered to give results regarding serious complications, hospitalization and mortality", as did a 2014 Cochrane Review. The Cochrane Review further recommended: "On the basis of the findings of this review, clinicians and healthcare policy-makers should urgently revise current recommendations for use of the neuraminidase inhibitors (NIs) for individuals with influenza." That is not utilizing NIs for prevention or treatment "Based on these findings there appears to be no evidence for patients, clinicians or policy-makers to use these drugs to prevent serious outcomes, both in annual influenza and pandemic influenza outbreaks."
The CDC, ECDC, Public Health England, Infectious Disease Society of America, the AAP, and Roche (the originator) reject the conclusions of the Cochrane Review, arguing in part that the analysis inappropriately forms conclusions about outcomes in people who are seriously ill based on results obtained primarily in healthy populations, and that the analysis inappropriately included results from people not infected with influenza. The EMA did not change its labeling of the drug in response to the Cochrane study.
A 2014 review recommended that all people admitted to intensive care units during influenza outbreaks with a diagnosis of community-acquired pneumonia receive oseltamivir until the absence of influenza infection is established by polymerase chain reaction (PCR) testing.
A 2015 systematic review and meta-analysis found oseltamivir effective at treating the symptoms of influenza, reducing the length of hospitalization, and reducing the risk of otitis media. The same review found that oseltamivir did not significantly increase the risk of adverse events. A 2016 systematic review found that oseltamivir slightly reduced the time it takes for the symptoms of influenza to be alleviated, and that it also increased the risk of "nausea, vomiting, [and] psychiatric events in adults and vomiting in children." The decrease in duration of sickness was about 18 hours.
Otherwise healthy people
In those who are otherwise healthy the CDC states that antivirals may be considered within the first 48 hours. A German clinical practice guideline recommends against its use.
Two 2013 meta-analyses have concluded that benefits in those who are otherwise healthy do not outweigh its risks. When the analysis was restricted to people with confirmed infection, the same 2014 Cochrane Review (see above) found unclear evidence of change in the risk of complications such as pneumonia, while three other reviews found a decreased risk. Together, published studies suggest that oseltamivir reduces the duration of symptoms by 0.5–1.0 day. Any benefit of treatment must be balanced against side effects, which include psychiatric symptoms and increased rates of vomiting.
The 2014 Cochrane Collaboration review concluded that oseltamivir did not affect the need for hospitalizations, and that there is no proof of reduction of complications of influenza (such as pneumonia) because of a lack of diagnostic definitions, or reduction of the spread of the virus. There was also evidence that suggested that oseltamivir prevented some people from producing sufficient numbers of their own antibodies to fight infection. The authors recommended that guidance should be revised to take account of the evidence of small benefit and increased risk of harms.
The US Centers for Disease Control and Prevention (CDC), the European Centre for Disease Prevention and Control (ECDC), the Public Health England (PHE), the Infectious Disease Society of America (IDSA), the American Academy of Pediatrics (AAP), and Roche (the originator) rejected the recommendations of the 2014 Cochrane Review to urgently change treatment guidelines and drug labels.
Prevention
, the CDC does not recommend to use oseltamivir generally for prevention due to concerns that widespread use will encourage resistance development. They recommend that it be considered in those at high risk, who have been exposed to influenza within 48 hours and have not received or only recently been vaccinated. They recommended it during outbreaks in long term care facilities and in those who are significantly immunosuppressed.
, reviews concluded that when oseltamivir is used preventatively it decreases the risk of exposed people developing symptomatic disease. A systematic review of systematic reviews found low to moderate evidence that it decreases the risk of getting symptomatic influenza by 1 to 12% (a relative decrease of 64 to 92%). It recommended against its use in healthy, low-risk persons due to cost, the risk of resistance development, and side effects and concluded it might be useful for prevention in unvaccinated high risk persons.
Side effects
Common adverse drug reactions (ADRs) associated with oseltamivir therapy (occurring in over 1 percent of people) include nausea and vomiting. In adults, oseltamivir increased the risk of nausea for which the number needed to harm was 28 and for vomiting was 22. So, for every 22 adult people on oseltamivir one experienced vomiting. In the treatment of children, oseltamivir also induced vomiting. The number needed to harm was 19. So, for every 19 children on oseltamivir one experienced vomiting. In prevention there were more headaches, kidney, and psychiatric events. Oseltamivir's effect on the heart is unclear: it may reduce cardiac symptoms, but may also induce serious arrhythmias.
Postmarketing reports include liver inflammation and elevated liver enzymes, rash, allergic reactions including anaphylaxis, toxic epidermal necrolysis, abnormal heart rhythms, seizure, confusion, aggravation of diabetes, and haemorrhagic colitis and Stevens–Johnson syndrome. Neuropsychiatric symptoms and body pain, followed by skin problems, were identified as the side effects most significantly reducing patient satisfaction on a drug review platform.
The US and EU package inserts for oseltamivir contain a warning of psychiatric effects observed in post-marketing surveillance. The frequency of these appears to be low and a causative role for oseltamivir has not been established. The 2014 Cochrane Review found a dose-response effect on psychiatric events. In trials of prevention in adults one person was harmed for every 94 treated. Neither of the two most cited published treatment trials of oseltamivir reported any drug-attributable serious adverse events.
It is pregnancy category B in Australia, meaning that it has been taken by a small number of women without signs of problems and in animal studies it looks safe. Dose adjustment may be needed in those with kidney problems.
Mechanism of action
Oseltamivir is a neuraminidase inhibitor, a competitive inhibitor of influenza's neuraminidase enzyme. The enzyme cleaves the sialic acid which is found on glycoproteins on the surface of human cells that helps new virions to exit the cell, preventing new viral particles from being released.
Resistance
The vast majority of mutations conferring resistance are single amino acid residue substitutions (His274Tyr in N1) in the neuraminidase enzyme. A 2011 meta-analysis of 15 studies found a pooled incidence rate for oseltamivir resistance of 2.6%. Subgroup analyses detected higher rates among influenza A patients, especially the H1N1 subtype. It was found that a substantial number of patients might become oseltamivir-resistant as a result of oseltamivir use, and that oseltamivir resistance might be significantly associated with pneumonia. In severely immunocompromised patients there were reports of prolonged shedding of oseltamivir- (or zanamivir)-resistant virus, even after oseltamivir treatment was stopped.
H1N1 flu or "swine flu"
, the World Health Organization (WHO) reported 314 samples of the prevalent 2009 pandemic H1N1 flu tested worldwide showed resistance to oseltamivir.
The CDC found sporadic oseltamivir-resistant 2009 H1N1 virus infections had been identified, including with rare episodes of limited transmission, but the public health impact had been limited. Those sporadic cases of resistance were found in immunosuppressed patients during oseltamivir treatment and persons who developed illness while receiving oseltamivir chemoprophylaxis.
During 2011, a new influenza A(H1N1)2009 variant with mildly reduced oseltamivir (and zanamivir) sensitivity was detected in more than 10% of community specimens in Singapore and more than 30% of samples from northern Australia.
While there is concern that antiviral resistance may develop in people with haematologic malignancies due to their inability to reduce viral loads and several surveillance studies found oseltamivir-resistant pH1N1 after administration of oseltamivir in those people, , widespread transmission of oseltamivir-resistant pH1N1 has not occurred.
During the 2007–08 flu season, the US CDC found 10.9% of H1N1 samples (n=1,020) to be resistant. In the 2008–09 season, the proportion of resistant H1N1 increased to 99.4%, while no other seasonal strains (H3N2, B) showed resistance.
Seasonal flu
From 2009 to 2014, oseltamivir resistance was very low in seasonal flu. In the 2010–11 flu season, 99.1% of H1N1, 99.8% of H3N, and 100% of Influenza B remained oseltamivir susceptible in the US. In January 2012, the US and European CDCs reported all seasonal flu samples tested since October 2011 to be oseltamivir susceptible. In the 2013–14 season only 1% of 2009 H1N1 viruses showed oseltamivir resistance. No other influenza viruses were resistant to oseltamivir.
H3N2
Three studies have found resistance in 0%, 3.3%, and 18% of subjects. In the study with the 18% resistance rate, the subjects were children, many of whom had not been previously exposed to influenza virus and therefore had a weakened immune response; the results suggest that higher and earlier dosing may be necessary in such populations.
Influenza B
In 2007, Japanese investigators detected neuraminidase-resistant influenza B virus strains in individuals not treated with these drugs. The prevalence was 1.7%. According to the CDC, , transmission of oseltamivir-resistant influenza B virus strains—from persons treated with the drug—is rare.
H5N1 avian influenza "bird flu"
, H274Y and N294S mutations that confer resistance to oseltamivir have been identified in a few H5N1 isolates from infected patients treated with oseltamivir, and have emerged spontaneously in Egypt.
H7N9 avian influenza
, two of 14 adults infected with A(H7N9) and treated with oseltamivir developed oseltamivir-resistant virus with the Arg292Lys mutation.
Pharmacokinetics
Its oral bioavailability is over 80% and is extensively metabolised to its active form upon first-pass through the liver. It has a volume of distribution of 23–26 litres. Its half-life is about 1–3 hours and its active carboxylate metabolite has a half-life of 6–10 hours. More than 90% of the oral dose is eliminated in the urine as the active metabolite.
History
Oseltamivir was discovered by scientists at Gilead Sciences using shikimic acid as a starting point for synthesis; shikimic acid was originally available only as an extract of Chinese star anise; but by 2006, 30% of the supply was manufactured recombinantly in E. coli. Gilead exclusively licensed their relevant patents to Roche in 1996. The drug has not been patent-protected in Thailand, the Philippines, Indonesia, and several other countries.
In 1999, the FDA approved oseltamivir phosphate for the treatment of influenza in adults based on two double-blind, randomized, placebo-controlled clinical trials. In June 2002, the European Medicines Agency (EMA) approved oseltamivir phosphate for prophylaxis and treatment of influenza. In 2003, a pooled analysis of ten randomised clinical trials concluded that oseltamivir reduced the risk of lower respiratory tract infections resulting in antibiotic use and hospital admissions in adults.
Oseltamivir (as Tamiflu) was widely used during the H5N1 avian influenza epidemic in Southeast Asia in 2005. In response to the epidemic, various governments – including those of the United Kingdom, Canada, Israel, United States, and Australia – stockpiled quantities of oseltamivir in preparation for a possible pandemic and there were worldwide shortages of the drug, driven by the high demand for stockpiling. In November 2005, US President George W. Bush requested that Congress fund US$1 billion for the production and stockpile of oseltamivir, after Congress had already approved $1.8 billion for military use of the drug. Defense Secretary Donald Rumsfeld, who was a past chairman of Gilead Sciences, recused himself from all government decisions regarding the drug.
In 2006, a Cochrane Review (since withdrawn) raised controversy by concluding that oseltamivir should not be used during routine seasonal influenza because of its low effectiveness.
In December 2008, the Indian drug company Cipla won a case in India's court system allowing it to manufacture a cheaper generic version of Tamiflu, called Antiflu. In May 2009, Cipla won approval from the World Health Organization (WHO) certifying that its drug Antiflu was as effective as Tamiflu, and Antiflu is included in the WHO list of prequalified medicinal products.
In 2009, a new A/H1N1 influenza virus was discovered to be spreading in North America. In June 2009, the WHO declared the A/H1N1 influenza a pandemic. The National Institute for Health and Care Excellence (NICE), the CDC, the WHO, and the ECDC maintained their recommendation to use oseltamivir.
From 2010 to 2012, Cochrane requested Roche's full clinical study reports of their trials, which they did not provide. In 2011, a freedom of information request to the European Medicines Agency (EMA) provided Cochrane with reports from 16 Roche oseltamivir trials. In 2012, the Cochrane team published an interim review based on those reports. In 2013, Roche released 74 full clinical study reports of oseltamivir trials after GSK released the data on zanamivir studies. In 2014, Cochrane published an updated review based solely on full clinical study reports and regulatory documents. In 2016, Roche's oseltamivir patents began to expire.
Veterinary use
There have been reports of oseltamivir reducing disease severity and hospitalization time in canine parvovirus infection. The drug may limit the ability of the virus to invade the crypt cells of the small intestine and decrease gastrointestinal bacterial colonization and toxin production.
Research
Oseltamivir has been deemed ineffective at treating COVID-19, consistent with the SARS-CoV-2 virus lacking influenza's neuraminidase enzyme.
| Biology and health sciences | Antiviral drugs | Health |
957233 | https://en.wikipedia.org/wiki/Mil%20Mi-8 | Mil Mi-8 | The Mil Mi-8 (, NATO reporting name: Hip) is a medium twin-turbine helicopter, originally designed by the Soviet Central Aerohydrodynamic Institute (TsAGI) in the 1960s and introduced into the Soviet Air Force in 1968. Russian production of the aircraft model still continues as of 2024. In addition to its most common role as a transport helicopter, the Mi-8 is also used as an airborne command post, armed gunship, and reconnaissance platform.
The Mi-8 is the world's most-produced helicopter, with over 17,000 units used by over 50 countries. As of 2015, when combined with the related Mil Mi-17, the two helicopters are the third most common operational military aircraft in the world.
Design and development
Mikhail Mil originally approached the Soviet government with a proposal to design an all-new two-engined turbine helicopter in 1959 after the success of the Mil Mi-4 and the emergence and effectiveness of turbines used in the Mil Mi-6. After design and development, the Mi-8 was subsequently introduced into the Soviet Air Force in 1967.
The Soviet military originally argued against a new helicopter, as they were content with the current Mil Mi-4. To counter this, Mikhail Mil proposed that the new helicopter was more of an update to new turbine engines rather than an entirely new helicopter, which persuaded the council of ministers to proceed with production. Due to the position of the engine, this enabled Mikhail Mil to justify redesigning the entire front half of the aircraft around the single engine.
The prototype, which was named V-8, was designed in 1958 and based on the Mil Mi-4 with a larger cabin. Powered by an AI-24 2,010 kW (2,700 shp) Soloviev turboshaft engine, the single engined V-8 prototype had its maiden flight in June 1961 and was first shown on Soviet Aviation Day parade (Tushino Air Parade) in July 1961.
During an official visit to the United States in September 1959, Nikita Khrushchev took a flight in the S-58 presidential helicopter for the first time and was reportedly extremely impressed. On Khrushchev's return, he ordered the creation of a similar helicopter, which was to be ready for the return visit by the American president, to save face. A luxury version of the Mi-4 was quickly created and Khrushchev took an inspection flight, during which Mikhail Mil proposed that his helicopter in development was more suitable. However, it would be necessary to have a second engine for reliability.
This gave Mikhail Mil the power under the orders of Khrushchev to build the original two-engined helicopter, which for the first time in Soviet history would need purpose-built turbine engines, rather than those adapted from fixed wing aircraft (as in the Mil Mi-6 and the first prototype V-8) and an entirely new main rotor gear box that would be designed in-house for the first time. In May 1960, the order was given for Mikhail Mil to create his twin engine helicopter. The Sergei Isotov Design Bureau accepted the task of creating the engines.
The second prototype (still equipped with the one turbine engine as the Isotov engines were still under development) flew in September 1961.
Two months after the engines were completed by Isotov, the third prototype designated V-8A equipped with two 1,120 kW (1,500 shp) Isotov TV2 engines, made its first flight piloted by Nikolai Ilyushin on 2 August 1962, marking the first flight of any Soviet helicopter to fly with purpose built gas turbine engines. The aircraft completed its factory based testing in February 1963.
The fourth prototype was designed as a VIP transport, with the rotor changed from four blades to five blades in 1963 to reduce vibration, the cockpit doors replaced by blister perspex slides and a sliding door added to the cabin.
The fifth and final prototype was a mass production prototype for the passenger market. In November 1964, all joint testing had been completed and the Soviet government began mass production. Production started in the Kazan Production Plant, with the first aircraft completed by the end of 1965.
The Soviet military originally showed little interest in the Mi-8 until the Bell UH-1's involvement in the Vietnam War became widely publicised as a great asset to the United States, allowing troops to move swiftly in and out of a battlefield and throughout the country. It was only then that the Soviet military rushed a troop-carrying variant of the Mil Mi-8 into production. By 1967, it had been introduced into the Soviet Air Force as the Mi-8.
There are numerous variants, including the Mi-8T, which, in addition to carrying 24 troops, is armed with rockets and anti-tank guided missiles. The Mil Mi-17 export version is employed by around 20 countries; its equivalent in Russian service in the Mi-8M series. The only visible differences between the Mi-8 and Mi-17 are A) the position of the tail rotor (Mi-8 right side, Mi-17 left side), B) the shape of the exhausts (Mi-8 circular, Mi-17 oval), and C) Dust shields in front of engine air intakes for the Mi-17. Also Mi-17 has some improved armour plating for its crew. The naval Mil Mi-14 version is also derived from the Mi-8.
The Mi-8 is constantly improving and the newest version still remains in production in 2024. However the second generation of the Mi-8 was changed to a tractor-tail rotor configuration as this configuration has increased yaw authority from the upwards advancing tail rotor blades into the downwash. The increase of the airspeed flowing over the rotor blades increases overall tail rotor effectiveness and yaw authority, whereas with the 'Pusher' tail rotor configuration the advancing rotor blade moves downwards. This decreases the airspeed across the rotor blade, reducing its overall effective yaw authority.
Operational history
Finland
The Finnish Defence Forces and the Finnish Border Guard began using Mi-8s in the 1970s, with the Finnish Air Force receiving its first, serialed HS-2, on 28 May 1973, and the second, HS-1, on 31 May 1973. Six Mi-8Ts were obtained at first, followed by further two Mi-8Ts and two Mi-8Ps. Three of the helicopters were handed over to the Border Guard Wing. One of these was lost after sinking through ice during a landing in April 1982. It was soon replaced by a new Mi-8.
After their Border Guard service, the helicopters were transferred to the civil register, but shortly thereafter to the Finnish Air Force. In 1997 it was decided that all helicopters, including the remaining five Mi-8Ts and two Mi-8Ps, should be transferred to the Army Wing at Utti. All Mi-8s have now been retired. One Mi-8 is on display at the Finnish Aviation Museum in Vantaa, and one is at the Päijänne Tavastia Aviation Museum in Asikkala, near Lahti. The two final Mi-8Ts were given to Hungary in August 2011 with all the remaining spare parts.
Georgia
The Georgian air force started operating Mi-8 and Mi-17 helicopters from 1991 onwards. During the War in Abkhazia (1992–1993) Mi-8 helicopters were used by both sides. Several were shot down, the first being a Georgian civilian Mi-8T which was destroyed in Sukhumi by an RPG-7. On 14 December 1992, a Russian Air Force Mi-8T was shot down by a SA-14 missile near Lata.
On another occasions Abkhaz Mi-8MTVs were shot down by Georgian forces, by SA-14 in one case and by RPG-18 in a second case, both during 1993. In the final case, Georgian Mi-8MTV carrying civilian refugees was shot down, killing 25 people. Georgian Air Force and Police currently operate about 20 Mi-8T/MTVs.
Iraq
Mi-8s were employed by the former Iraqi Army Aviation and Iraqi Air Force under Saddam Hussein. In the Iran–Iraq War of the 1980s, there were air-to-air combat between Iraqi and Iranian Army Aviation helicopters, including between Iranian Bell AH-1J Cobras and Iraqi Mi-8s.
South Sudan
On 21 December 2012, a Nizhnevartovskavia owned Mi-8 working for the United Nations Mission in South Sudan (UNMISS) was shot down and crashed near Likuangole in the South Sudanese state of Jonglei during the South Sudan internal conflict. All four Russian crewmembers on board were killed, and after some initial confusion, a UN spokesman said that the South Sudanese army confirmed on 22 December that it mistakenly fired at the helicopter.
On 26 August 2014, a UTair Aviation owned Mi-8 working for the United Nations crashed as it approached a landing airstrip near Bentiu. Three of the Russian crew members died and one was injured. Rebel commander Peter Gadet claimed that his forces brought it down using a rocket-propelled grenade.
Soviet Union
The Mi-8 family of helicopters became the main Soviet (and later Russian) helicopter, covering a large range of roles in both peace time and war time. Large fleets of Mi-8 and its derivatives were employed by both military and civil operators.
Large numbers of Mi-8 family helicopters were used during the Soviet–Afghan War during the 1980s. Its rugged construction allowed easier in-theater operations and maintenance. A large number of Mi-8s were lost with several shot down by enemy fire, with the Mi-8 and its derivatives being the main aircraft model lost by the Soviet Union in Afghanistan.
Between April and May 1986, Mi-8s were used in large numbers to drop radiation-absorbing materials into the No. 4 reactor of Chernobyl Nuclear Power Plant after the Chernobyl disaster, and the fire was extinguished by the combined effort of helicopters dropping over 5000 metric tons of sand, lead, clay, and neutron-absorbing boron onto the burning reactor and injecting liquid nitrogen into it. Most of the helicopters were severely irradiated and abandoned in a giant junkyard, the so-called "machines cemetery" near Chernobyl, with several disappearing from the site in later years. During the initial operation, one crashed near the power plant after hitting a construction crane cable with all the crew of four being killed in the crash. It is now known that virtually none of the neutron absorbers reached the core.
Ukraine
On 16 August 2013, the Ministry of Defense of Ukraine reported that one of its Mi-8MSB had set a world altitude record of at the Kirovske military airfield on 15 August.
The Ukrainian Armed Forces used Mi-8MSB along with Mi-24s in operations against separatists in Eastern Ukraine during the Russo-Ukrainian War. On 29 May 2014, a Ukrainian National Guard Mi-8 was brought down by Russian separatist forces in Donbas using a MANPADS near Slavyansk with 12 personnel, including an Army general, killed and one seriously injured. On 24 June 2014, a Ukrainian National Guard Mi-8 was shot down by separatist forces again using a MANPADS near Slavyansk with nine personnel killed. Ukrainian forces used Mi-8 helicopters to resupply forces during the Siege of Mariupol at Azovstal iron and steel works and bring in additional reinforcements for the Azov Regiment. Some 16 Mi-8s were used a number of times, two of which were shot down. Russia claimed on 5 April that it shot down two Ukrainian Mi-8s that it said were being used to evacuate commanders of the Azov Regiment. In late August 2023, it was reported that a Russian defector named Maksym Kuzminov handed over a Mi-8AMTSh to the Ukrainian forces in coordination with Ukrainian Intelligence agents.
On 16 October 2023, Ukrainian Colonel General Oleksandr Syrskyi said that the 25th Separate Airborne Brigade had shot down a Mi-8 without supplying the location. According to the Oryx database 21 Mi-8 helicopters have been shot down by Ukraine thus far during the war.
On 31 July 2024, a Russian Mi-8 helicopter was shot down by over occupied Donetsk, by Ukrainian FPV drones. The first time a helicopter in combat was destroyed by a drone. The Mi-8 was believed to have been attacked on the ground either during landing or take off.
On 31 December 2024, a Russian Mi-8 was shot down by a MAGURA V5 sea drone armed with R-73 Sea Dragon missiles near Cape Tarkhankut, Crimea, while a second helicopter was damaged but managed to return to base, according to the Main Directorate of Intelligence.
United States
During the initial stages of Operation Enduring Freedom, Mi-17s and Mi-8s were extensively used by the CIA and US Special Forces to assist the Northern Alliance in their fight against the Taliban.
A number of Mi-8s and Mi-17s are used by US government agencies as of 2022.
Yugoslavia
The Yugoslav Air Force took delivery of 24 Mi-8T (Hip C) transport helicopters between May 1968 and May 1969 to equip two squadrons of the newly formed 119th transport regiment from Niš military airport, each squadron with 12 helicopters. Subsequently, from 1973 to the early 1980s, Yugoslavia purchased more Mi-8T helicopters to re-equip two squadrons of 111th regiment from Pleso military airport near Zagreb and the 790th squadron from Divulje military airport near Split, which was under the command of the Yugoslav Navy. In total, the Yugoslav Air Force received 92 Mi-8Ts, designated by the Yugoslav People's Army as the HT-40, while local modification of several helicopters into electronic warfare variants produced the HT-40E. Some 40 helicopters were equipped for firefighting operations.
The Yugoslav Mi-8s' first combat operations were transport of Yugoslav People's Army troops and federal police forces to border crossings in Slovenia on 27 June 1991 during the Ten-Day War. The members of Slovenian Territorial Defence fired Strela 2 MANPAD, and shot one helicopter down, killing all crew and passengers.
During combat in the winter of 1991 in the Croatian war and in the spring of 1992 in the Bosnian War, the Yugoslav People's Army used the Mi-8 fleet for the evacuation of injured personnel, transport of cargo and search and rescue for the crews of aircraft forced down. As most flights were made behind the front, the Croatian forces were able to down just one helicopter, which was hit by small arms fire near Slavonski Brod on 4 October 1991.
After Bosnian Serbs declared their state in the spring of 1992, some former Yugoslav Air Force Mi-8s continued service with the Republika Srpska armed forces. The inventory of the 82nd mixed helicopter squadron, of the 92nd aviation brigade of the Army of Republika Srpska comprised 12 Mi-8T helicopters, which continued in service until Operation Koridor. During that period, the Republika Srpska Air Force lost three Mi-8 helicopters to enemy fire. Three helicopters painted in a blue and white colour scheme flew in the first part of 56th helicopter squadron of the Krajina Milicija, using Udbina military airport in Lika as their main base. The Republika Srpska Air Force continued to operate nine helicopters, albeit suffering problems with maintenance and spare parts, until it was formally disbanded in 2006.
On the other side, Mi-8 helicopters were also used as main air transport. The Croatian National Guard obtained its first on 23 September 1991, near Petrinja, when a Yugoslav Air Force Mi-8 made an emergency landing after being damaged by small-arms fire. A further 6 Mi-8T and 18 Mi-8MTV-1 helicopters were bought from ex-Warsaw Pact countries during the war, with 16 being used in active service, and remaining were used as source for spare parts. The remaining Mi-8Ts were retired from service in the Croatian Air Force after the war, while the Mi-8MTVs continued their service in 20th Transport Helicopter Squadron and 28th Transport Helicopter Squadron. The latter has been re-equipped with new Mi-171Sh helicopters bought from Russia.
The Army of the Republic of Bosnia and Herzegovina secretly obtained Mi-8T, Mi-8MTV and Mi-17 helicopters from various sources. Two helicopters were shot down by Serb air defenses, one around Žepa, while one Mi-17 was shot down by 2K12 Kub M, killing the Bosnian Foreign Affairs Minister Irfan Ljubijankić, a few other politicians, and the helicopter's Ukrainian crew. A few Croatian Mi-8MTVs secretly supported Croatian Defence Council operations in Herceg Bosna. After the war, the Army of the Federation of Bosnia and Herzegovina operated the remaining five Mi-8MTVs and one Mi-8T in the Air Force and Air Defense Brigade of Armed Forces of Bosnia and Herzegovina.
The North Macedonian Air Force bought two Mi-8MT helicopters in 2001 from Ukraine. They fly in the Transport Helicopter Squadron (ex 301. Transport Helicopter Squadron). One crashed, killing all 8 passengers and 3 crew members in an accident in January 2008.
During the Kosovo War of 1998 and 1999, the Federal Yugoslav Air Force used Mi-8s for transport of personnel and material to forces in otherwise-inaccessible mountain areas. Evacuation of injured personnel also occurred during the 1999 NATO bombing of Yugoslavia, flying at low altitude to avoid detection by NATO aircraft. In 1999, Yugoslav Mi-8s shot down at least one US Army Hunter UAV with the door gunner's 7.62 mm machine gun. Two Mi-17V helicopters secretly operated by the Special Operations Unit post-1997 were also active during the Kosovar conflict. After the unit disbanded in 2003, the helicopters were transferred to Serbia and Montenegro's air force.
As of mid-2020, the Serbian Air Force, the successor of the Federal Yugoslav Air Force, operates a small amount of Mi-8T which are now being replaced by Mi-17 helicopters. There are 13 Mi-17 in the Serbian air force currently. They are in the 138th Mixed-Transport-Aviation Squadron of 204th Air Base and 119th Combined-Arms Helicopter Squadron (ex 199th regiment) of 98th Air Base.
Others
Canada – After Canada committed combat forces to fight the Taliban in Afghanistan, they realized their mobility depended on borrowed helicopter airlift. In 2007, the Minister of National Defence Peter MacKay announced the lease of 6 to 8 Mi-8s, particularly Kazan Helicopters Mi-17-V5s, until the introduction of 6 interim CH-47Ds in 2008 and later delivery of 15 new-build CH-47Fs in 2013 by the RCAF.
Poland – On 4 December 2003, a Polish Mi-8 crashed near Piaseczno while carrying Prime Minister Leszek Miller, ten other passengers and four crewmen. There were no fatalities. The cause of the accident was the icing of the engines. The pilot was accused of causing the crash, but he was found not guilty.
Syria – During the Yom Kippur War of October 1973, Syria landed special forces troops behind Israel Defense Forces lines on the Golan Heights at Mt. Hermon, Tel Fares, Vaset, Nafach and Ein Zivan – A Dalve.
Yemen—On 19 November 2023, Houthi rebels utilized a captured Mi-17 helicopter to conduct an air assault boarding and seizure of the Japanese owned cargo ship Galaxy Leader.
Vietnam- On 22 November 1992 a Vietnamese Mi-8 was sent from Hanoi carrying rescue workers for Flight 474, but it crashed near Ô Kha mountain on the same day. All seven people aboard were killed
Variants
Prototypes/experimental/low production rate
V-8 (NATO – Hip-A)
The original single-engined prototype.
V-8A
A twin-engined prototype, featuring TV2-117 turboshaft engines, the prototype underwent further modifications during its life.
V-8AT
Prototype of the Mi-8T utility version.
Mi-8 (NATO – Hip-B)
Twin-engined prototype.
Mi-8TG
Conversion to operate on LPG gas.
Mi-18
Prototype design, a modification of the existing Mil Mi-8. Two Mi-8s were extended by 0.9 meters (3 ft), the landing gear made retractable, and a sliding door added to the starboard side of the fuselage. The Mi-18s were used in the Soviet invasion of Afghanistan, and later used as static training airframes for pilots of the Mi-8/17.
Basic military transport/airframe
Mi-8T (NATO – Hip-C)
First mass production utility transport version, it can carry four UV-16-57 unguided rocket pods, (with S-5 rockets), mounted to four hardpoints on two outrigger pylons, and is armed with one or two side-mounted PK machine guns.
Mi-8TV
Armed version of the Mi-8T.
Mi-8TVK (NATO – Hip-E, a.k.a. Mi-8TB)
Version used as a gunship or direct air support platform. Airframe modifications add 2x external hard points for a total of 6, and mount a flexible 12.7 mm (0.5-inch) KV-4 machine gun in the nose. Armament of 57 mm S-5 rockets, six UV-32-57 rocket pods, 551-lb (250-kg) bombs, or four AT-2 Swatter ATGMs.
Mi-8TBK (NATO – Hip-F)
Armed export version, fitted with six launch rails to carry and fire Malyutka missiles.
Command and electronic warfare
Mi-8IV (NATO – Hip-G, a.k.a. Mi-9)
Airborne command post version fitted with "Ivolga" system, characterized by antennas, and Doppler radar on tail boom.
Mi-8PP (NATO – Hip-K)
Airborne jamming platform with "Polye" (field) system. From 1980, the type was fitted with the new "Akatsiya" system and redesignated the Mi-8PPA. It is characterized by six X-shaped antennas on each side of the aft fuselage. Built to escort troop-carrying versions of this helicopter, and disrupt potentially nearby SPAAG radars, such as those of the Flakpanzer Gepard.
Mi-8PD
Polish airborne command post version.
Mi-8SMV (NATO – Hip-J)
Airborne jamming platform with "Smalta-V" system, characterized by two small boxes on each side of the fuselage. Used for protection of ground attack aircraft against enemy air defenses.
Mi-8VKP (NATO – Hip-D, a.k.a. Mi-8VzPU)
Airborne communications platform with rectangular communication canisters mounted on weapons racks and with two frame-type aerials above the rear fuselage.
Other military
Mi-8AD
Minelaying version with four VSM-1 dispensers.
Mi-8AV
Minelaying version with VMR-1 or −2 system for 64 or 200 anti-tank mines.
Mi-8BT
Mine-clearing version.
Mi-8MB "Bissektrisa"
Military ambulance version.
Mi-8R (a.k.a. Mi-8GR)
Tactical reconnaissance version with Elint system "Grebeshok-5".
Mi-8K
Artillery observation, reconnaissance version.
Mi-8SMT
Military staff transport version, fitted with improved radio equipment R-832 and R-111.
Mi-8SKA
Photo-reconnaissance version.
Mi-8SP
Spacecraft tracking and recovery version.
Mi-8T(K)
Photo-reconnaissance version.
Mi-8TZ
Fuel transport tanker version.
Mi-8TB
The Mi-8TB was developed in the GDR and specially adapted to the military needs there. It was equipped with various missile and bomb systems, including S-5 missiles and FAB-500 bombs, which could be used to attack ground targets. These adjustments made them a type of “transport bomber,” which explained the “TB” designation.
Mi-8MTYu
Only one was built and used by the Ukrainian Air Force, based at AB "Kirovske". Intended for detection of re-entry vehicles, and small surface targets. In the nose radar antenna.
Mi-8MSB
Modernized passenger-transport version for civil aviation.
Mi-8MSB-V
Modernized multipurpose helicopter for the Ukrainian Armed Forces.
Civil
Mi-8T (NATO – Hip-C)
Civilian and military utility transport version, with accommodation for 24 passengers, fitted with tip-up seats along the cabin walls, circular cabin windows and large rear clamshell doors with a sloping hinge line. The Mi-8T is powered by two Klimov TV2-117A turboshaft engines, giving the helicopter a maximum speed of at sea level.
Mi-8P
Civilian passenger transport version, with accommodation for between 28 and 32 passengers, fitted with square cabin windows, small rear clamshell doors with a vertical hinge line and a horizontally split rear airstair door in between; powered by two Klimov TV2-117A turboshaft engines.
Mi-8S "Salon"
Civilian VIP transport version, with accommodation for between 9 and 11 passengers, equipped with a galley and toilet.
Mi-8MPS
Search and rescue version (operated usually in Malaysia for Fire and Rescue Department services).
Mi-8MA
Polar exploration version for use in the Arctic.
Mi-8MT
Flying crane version.
Mi-8AT
Civilian transport version, fitted with two improved TV2-117AG turboshaft engines.
Mi-8ATS
Agricultural version, fitted with a hopper and spray bars.
Mi-8TL
Air accident investigation version.
Mi-8TM
Upgraded transport version, fitted with a weather radar.
Mi-8TS
Hot and high desert version.
Mi-8VIP
Deluxe VIP transport version, with accommodation for between 7 and 9 passengers.
Mi-8PA
Modified version for Japanese regulations. One only was built, in 1980. It was used by Aero Asahi for heavy material transport in a mountainous region. It was retired in 1993 and later moved to the Tokorozawa Aviation Museum.
Accidents and incidents
On 1 November 1974, a Mil Mi-8T helicopter collided with a Antonov An-2 in Surgut, Soviet Union, killing all 24 people on the helicopter and all 14 people on the Antonov An-2.
On 18 September 1981, a Mil Mi-8T helicopter collided with a Yakovlev Yak-40 east of Zheleznogorsk-Ilimsky, Russia, killing all seven crew on the helicopter and all thirty-three people on the Yak-40.
On 17 September 2001, a Mil Mi-8 helicopter was shotdown by a Chechen group in Grozny, Chechnya, Russia, killing all 13 people on board.
On 27 January 2002, a Mil Mi-8 helicopter was shotdown in Shelkovskaya, Nadterechny District, Russia, killing all 14 people on board.
On 4 December 2003, a Mil Mi-8P helicopter crashed near Piaseczno, Poland, injuring 8 or 14 people on board, including Polish Prime Minister Leszek Miller.
On 11 September 2006, a Mil Mi-8 helicopter was shotdown by militant group Kataib al-Khoul near Vladikavkaz, Russia, killing 12 of the 16 people on board.
On 23 September 2006, a Mil Mi-8 MTV 1 helicopter crashed in Ghunsa, Nepal, killing all 24 people on board.
On 27 April 2007, a Mil Mi-8 helicopter crashed near Shatoy, Chechnya, Russia, killing all 20 people on board.
On 3 June 2007, a Mil Mi-8 helicopter crashed in Lungi, Sierra Leone, killing all 22 people on board.
On 7 April 2013, a Mil Mi-8PS helicopter crashed in Loreto, Peru, killing all 13 people on board.
On 2 July 2013, a Mil Mi-8 helicopter crashed near Deputatsky, Sakha Republic, Russia, killing 24 of the 28 people on board.
On 21 October 2016, a Mil Mi-8T helicopter crashed in Yamalo-Nenets Autonomous Okrug, Yamal Peninsula, Siberia, Russia, killing 19 of the 22 people on board.
On 12 August 2021, a Mil Mi-8T helicopter into Kurile Lake, Kamchatka Krai, Russia, killing eight people.
On 31 August 2024, a Mil Mi-8T helicopter crashed in Kamchatka Peninsula, Russia, killing all 22 occupants onboard.
On 28 September 2024, a Mil Mi-8 MTV-1 helicopter crashed in Waziristan, Pakistan, killing 6 occupants out of the 15 onboard.
Operators
(leased from SkyLink Aviation)
: 5
* (captured).
Former operators
Specifications (Mi-8MT)
| Technology | Specific aircraft_2 | null |
957422 | https://en.wikipedia.org/wiki/Keep | Keep | A keep is a type of fortified tower built within castles during the Middle Ages by European nobility. Scholars have debated the scope of the word keep, but usually consider it to refer to large towers in castles that were fortified residences, used as a refuge of last resort should the rest of the castle fall to an adversary. The first keeps were made of timber and formed a key part of the motte-and-bailey castles that emerged in Normandy and Anjou during the 10th century; the design spread to England, Portugal, south Italy and Sicily. As a result of the Norman Conquest of England in 1066, use spread into Wales during the second half of the 11th century and into Ireland in the 1170s. The Anglo-Normans and French rulers began to build stone keeps during the 10th and 11th centuries, including Norman keeps, with a square or rectangular design, and circular shell keeps. Stone keeps carried considerable political as well as military importance and could take a decade or more to build.
During the 12th century, new designs began to be introduced – in France, quatrefoil-shaped keeps were introduced, while in England polygonal towers were built. By the end of the century, French and English keep designs began to diverge: Philip II of France built a sequence of circular keeps as part of his bid to stamp his royal authority on his new territories, while in England castles were built without keeps. In Spain, keeps were increasingly incorporated into both Christian and Islamic castles, although in Germany tall fighting towers called bergfriede were preferred to keeps in the western fashion. In the second half of the 14th century, there was a resurgence in the building of keeps. In France, the keep at Vincennes near Paris began a fashion for tall, heavily machicolated designs, a trend adopted in Spain most prominently through the Valladolid school of Spanish castle design. Meanwhile, tower keeps in England became popular amongst the most wealthy nobles: these large keeps, each uniquely designed, formed part of the grandest castles built during the period.
In the 15th century, the protective function of keeps was compromised by improved artillery. For example, in 1464 during the Wars of the Roses, the keep of Bamburgh Castle on the Northumberland coast, previously considered to be impregnable, was defeated with bombards. By the 16th century, keeps were slowly falling out of fashion as fortifications and residences. Many were destroyed in civil wars between the 17th and 18th centuries or incorporated into gardens as an alternative to follies. During the 19th century, keeps became fashionable once again, and in England and France, a number were restored or redesigned by Gothic architects. Despite further damage to many French and Spanish keeps during the wars of the 20th century, keeps now form an important part of the tourist and heritage industry in Europe.
Etymology and historiography
Since the 16th century, the English word keep has commonly referred to large towers in castles. The word originates from around 1375 to 1376, coming from the Middle English term kype, meaning basket or cask, and was a term applied to the shell keep at Guînes, said to resemble a barrel. The term came to be used for other shell keeps by the 15th century. By the 17th century, the word keep lost its original reference to baskets or casks and was popularly assumed to have come from the Middle English word keep, meaning to hold or to protect.
Early on, the use of the word keep became associated with the idea of a tower in a castle that would serve both as a fortified, high-status private residence and a refuge of last resort. The issue was complicated by the building of fortified Renaissance towers in Italy called tenazza that were used as defences of last resort and were also named after the Italian for to hold or to keep. By the 19th century, Victorian historians incorrectly concluded that the etymology of the words "keep" and tenazza were linked and that all keeps had fulfilled this military function.
As a result of this evolution in meaning, the use of the term keep in historical analysis today can be problematic. Contemporary medieval writers used various terms for the buildings we would today call keeps. In Latin, they are variously described as turris, turris castri or magna turris – a tower, a castle tower, or a great tower. The 12th-century French came to term them a donjon, from the Latin dominarium "lordship", linking the keep and feudal authority. Similarly, medieval Spanish writers called the buildings torre del homenaje, or "tower of homage". In England, donjon turned into dungeon, which initially referred to a keep, rather than to a place of imprisonment.
While the term remains in common academic use, some academics prefer to use the term donjon, and most modern historians warn against using the term "keep" simplistically. The fortifications that we would today call keeps did not necessarily form part of a unified medieval style, nor were they all used in a similar fashion during the period.
History
Timber keeps (9th–12th centuries)
The earliest keeps were built as part of motte-and-bailey castles from the 10th century onwards – a combination of documentary and archaeological evidence places the first such castle, built at Vincy, in 979. These castles were initially built by the more powerful lords of Anjou in the late 10th and 11th centuries, in particular Fulk III and his son, Geoffrey II, who built a great number of them between 987 and 1060. William the Conqueror then introduced this form of castle into England when he invaded in 1066, and the design spread through south Wales as the Normans expanded up the valleys during the subsequent decades.
In a motte-and-bailey design, a castle would include a mound called a motte, usually artificially constructed by piling up turf and soil, and a bailey, a lower walled enclosure. A keep and a protective wall would usually be built on top of the motte. Some protective walls around a keep would be large enough to have a wall-walk around them, and the outer walls of the motte and the wall-walk could be strengthened by filling in the gap between the wooden walls with earth and stones, allowing it to carry more weight – this was called a garillum. Smaller mottes could only support simple towers with room for a few soldiers, whilst larger mottes could be equipped with a much grander keep. Many wooden keeps were designed with a bretasche, a square structure that overhung from the upper floors of the building, enabling better defences and a more sturdy structural design. These wooden keeps could be protected by skins and hides to prevent them from being easily set alight during a siege.
One contemporary account of these keeps comes from Jean de Colmieu around 1130, who described how the nobles of the Calais region would build "a mound of earth as high as they can and dig a ditch about it as wide and deep as possible. The space on top of the mound is enclosed by a palisade of very strong hewn logs, strengthened at intervals by as many towers as their means can provide. Inside the enclosure is a citadel, or keep, which commands the whole circuit of the defences. The entrance to the fortress is by means of a bridge, which, rising from the outer side of the moat and supported on posts as it ascends, reches to the top of the mound." At Durham Castle, contemporaries described how the keep arose from the "tumulus of rising earth" with a keep reaching "into thin air, strong within and without", a "stalwart house...glittering with beauty in every part". As well as having defensive value, keeps and mottes sent a powerful political message to the local population.
Wooden keeps could be quite extensive in size and, as Robert Higham and Philip Barker have noted, it was possible to build "...very tall and massive structures." As an example of what these keeps may have comprised, the early 12th-century chronicler Lambert of Ardres described the wooden keep on top of the motte at the castle of Ardres, where the "...first storey was on the surface of the ground, where were cellars and granaries, and great boxes, tuns, casks, and other domestic utensils. In the storey above were the dwelling and common living-rooms of the residents in which were the larders, the rooms of the bakers and butlers, and the great chamber in which the lord and his wife slept...In the upper storey of the house were garret rooms...In this storey also the watchmen and the servants appointed to keep the house took their sleep."
In the Holy Roman Empire, tall, free-standing, wooden (later stone), fighting towers called Bergfriede were commonly built by the 11th century, either as part of motte-and-bailey designs or, as part of Hohenburgen castles, with characteristic inner and outer courts. Bergfriede, which take their name from the German for a belfry, had similarities to keeps, but are usually distinguished from them on account of Bergfriede having a smaller area or footprint, usually being non-residential and being typically integrated into the outer defences of a castle, rather than being a safe refuge of last resort.
Early stone keeps (10th–12th centuries)
During the 10th century, a small number of stone keeps began to be built in France, such at the Château de Langeais: in the 11th century, their numbers increased as the style spread through Normandy across the rest of France and into England, South Italy and Sicily. Some existing motte-and-bailey castles were converted to stone, with the keep usually amongst the first parts to be upgraded, while in other cases new keeps were built from scratch in stone. These stone keeps were introduced into Ireland during the 1170s following the Norman occupation of the east of the country, where they were particularly popular amongst the new Anglo-Norman lords. Two broad types of design emerged across France and England during the period: four-sided stone keeps, known as Norman keeps or great keeps in English – a donjon carré or donjon roman in French – and circular shell keeps.
The reasons for the transition from timber to stone keeps are unclear, and the process was slow and uneven, taking many years to take effect across the various regions. Traditionally it was believed that stone keeps had been adopted because of the cruder nature of wooden buildings, the limited lifespan of wooden fortifications and their vulnerability to fire, but recent archaeological studies have shown that many wooden castles were as robust and as sophisticated as their stone equivalents. Some wooden keeps were not converted into stone for many years and were instead expanded in wood, such as at Hen Domen. Nonetheless, stone became increasingly popular as a building material for keeps for both military and symbolic reasons.
Stone keep construction required skilled craftsmen. Unlike timber and earthworks, which could be built using unfree labour or serfs, these craftsmen had to be paid and stone keeps were therefore expensive. They were also relatively slow to erect, due to the limitations of the lime mortar used during the period – a keep's walls could usually be raised by a maximum of only 12 feet (3.6 metres) a year; the keep at Scarborough was not atypical in taking ten years to build. The number of such keeps remained relatively low: in England, for example, although several early stone keeps had been built after the conquest, there were only somewhere between ten and fifteen in existence by 1100, and only around a hundred had been built by 1216.
Norman keeps had four sides, with the corners reinforced by pilaster buttresses; some keeps, particularly in Normandy and France, had a barlongue design, being rectangular in plan with their length twice their width, while others, particularly in England, formed a square. These keeps could be up to four storeys high, with the entrance placed on the first storey to prevent the door from being easily broken down; early French keeps had external stairs in wood, whilst later castles in both France and England built them in stone. In some cases the entrance stairs were protected by additional walls and a door, producing a forebuilding. The strength of the Norman design typically came from the thickness of the keep's walls: usually made of rag-stone, these could be up to 24 feet (7.3 metres) thick, immensely strong, and producing a steady temperature inside the building throughout summer and winter. The larger keeps were subdivided by an internal wall while the smaller versions had a single, slightly cramped chamber on each floor. Usually only the first floor would be vaulted in stone, with the higher storeys supported with timbers.
There has been extensive academic discussion of the extent to which Norman keeps were designed with a military or political function in mind, particularly in England. Earlier analyses of Norman keeps focused on their military design, and historians such as R. Brown Cathcart King proposed that square keeps were adopted because of their military superiority over timber keeps. Most of these Norman keeps were certainly extremely physically robust, even though the characteristic pilaster buttresses added little real architectural strength to the design. Many of the weaknesses inherent to their design were irrelevant during the early part of their history. The corners of square keeps were theoretically vulnerable to siege engines and galleried mining, but before the introduction of the trebuchet at the end of the 12th century, early artillery stood little practical chance of damaging the keeps, and galleried mining was rarely practised. Similarly, the corners of a square keep created dead space that defenders could not fire at, but missile fire in castle sieges was less important until the introduction of the crossbow in the middle of the 12th century, when arrowslits began to be introduced.
Nonetheless, many stone Norman keeps made considerable compromises to military utility. Norwich Castle, for example, included elaborate blind arcading on the outside of the building and appears to have had an entrance route designed for public ceremony, rather than for defence. The interior of the keep at Hedingham could certainly have hosted impressive ceremonies and events, but contained numerous flaws from a military perspective. Important early English and Welsh keeps such as the White Tower, Colchester, and Chepstow were all built in a distinctive Romanesque style, often reusing Roman materials and sites, and were almost certainly intended to impress and generate a political effect amongst local people. The political value of these keep designs, and the social prestige they lent to their builders, may help explain why they continued to be built in England into the late 12th century, beyond the point when military theory would have suggested that alternative designs were adopted.
The second early stone design, emerging from the 12th century onwards, was the shell keep, a donjon annulaire in French, which involved replacing the wooden keep on a motte, or the palisade on a ringwork, with a circular stone wall. Shell keeps were sometimes further protected by an additional low protective wall, called a chemise, around their base. Buildings could then be built around the inside of the shell, producing a small inner courtyard at the centre. The style was particularly popular in south-east England and across Normandy, although less so elsewhere. Restormel Castle is a classic example of this development, as is the later Launceston Castle; prominent Normandy and Low Country equivalents include Gisors and the Burcht van Leiden – these castles were amongst the most powerful fortifications of the period. Although the circular design held military advantages over one with square corners, as noted above these really mattered from only the end of the 12th century onwards; the major reason for adopting a shell keep design, in the 12th century at least, was the circular design of the original earthworks exploited to support the keep; indeed, some designs were less than circular in order to accommodate irregular mottes, such as that found at Windsor Castle.
Mid-medieval keeps (late 12th–14th centuries)
During the second half of the 12th century, a range of new keep designs began to appear across France and England, breaking the previous unity of the regional designs. The use of keeps in castles spread through Iberia, but some new castles never incorporated keeps in their designs. One traditional explanation for these developments emphasises the military utility of the new approaches, arguing, for example, that the curved surfaces of the new keeps helped to deflect attacks, or that they drew on lessons learnt during the Crusades from Islamic practices in the Levant. More recent historical analysis, however, has emphasised the political and social drivers that underlay these mid-medieval changes in keep design.
Through most of the 12th century, France was divided between the Capetian kings, ruling from the Île-de-France, and kings of England, who controlled Normandy and much of the west of France. Within the Capetian territories, early experimentation in new keep designs began at Houdan in 1120, where a circular keep was built with four round turrets; internally, however, the structure remained conventionally square. A few years later, Château d'Étampes adopted a quatrefoil design. These designs, however, remained isolated experiments.
In the 1190s, however, the struggle for power in France began to swing in favour of Philip II, culminating in the Capetian capture of Normandy in 1204. Philip II started to construct completely circular keeps, such as the Tour Jeanne d'Arc, with most built in his newly acquired territories. The first of Philip's new keeps was begun at the Louvre in 1190 and at least another twenty followed, all built to a consistent standard and cost. The architectural idea of circular keeps may have come from Catalonia, where circular towers in castles formed a local tradition, and probably carried some military advantages, but Philip's intention in building these new keeps in a fresh style was clearly political, an attempt to demonstrate his new power and authority over his extended territories. As historian Philippe Durand suggests, these keeps provided military security and were a physical representation of the renouveau capétien, or Capetian renewal.
Keep design in England began to change only towards the end of the 12th century, later than in France. Wooden keeps on mottes ceased to be built across most of England by the 1150s, although they continued to be erected in Wales and along the Welsh Marches. By the end of the 12th century, England and Ireland saw a handful of innovative angular or polygonal keeps built, including the keep at Orford Castle, with three rectangular, clasping towers built out from the high, circular central tower; the cross-shaped keep of Trim Castle and the famous polygonal design at Conisborough. Despite these new designs, square keeps remained popular across much of England and, as late as the 1170s, square Norman great keeps were being built at Newcastle. Circular keep designs similar to those in France really became popular in Britain in the Welsh Marches and Scotland for only a short period during the early 13th century.
As with the new keeps constructed in France, these Anglo-Norman designs were informed both by military thinking and by political drivers. The keep at Orford has been particularly extensively analysed in this regard, and although traditional explanations suggested that its unusual plan was the result of an experimental military design, more recent analysis concludes that the design was instead probably driven by political symbolism and the need for Henry to dominate the contested lands of East Anglia. The architecture would, for mid-12th century nobility, have summoned up images of King Arthur or Constantinople, then the idealised versions of royal and imperial power. Even formidable military designs such as that at Château Gaillard were built with political effect in mind. Gaillard was designed to reaffirm Angevin authority in a fiercely disputed conflict zone and the keep, although militarily impressive, contained only an anteroom and a royal audience chamber, and was built on soft chalk and without an internal well, both serious defects from a defensive perspective.
During most of the medieval period, Iberia was divided between Christian and Islamic kingdoms, neither of which traditionally built keeps, instead building watchtowers or mural towers. By the 12th century, however, the influence of France and the various military orders was encouraging the development of square keeps in Christian castles across the region, and by the second half of the century this practice was spread across into the Islamic kingdoms.
By contrast, the remainder of Europe saw stone towers being used in castles, but not in a way that fulfilled the range of functions seen in the western European keeps. In the Low Countries, it became popular for the local nobility to build stand-alone, square towers, but rarely as part of a wider castle. Similarly, square stone towers became popular in Venice, but these did not fulfil the same role as western keeps. In Germany, rectangular stone castles began to replace motte-and-bailey castles from the 12th century onwards. These designs included stone versions of the traditional Bergfriede, which still remained distinct from the domestic keeps used in more western parts of Europe, with the occasional notable exception, such as the large, residential Bergfried at Eltville Castle.
Several designs for new castles emerged that made keeps unnecessary. One such design was the concentric approach, involving exterior walls guarded with towers, and perhaps supported by further, concentric layered defenses: thus castles such as Framlingham never had a central keep. Military factors may well have driven this development: R. Brown, for example, suggests that designs with a separate keep and bailey system inherently lacked a co-ordinated and combined defensive system, and that once bailey walls were sophisticated enough, a keep became militarily unnecessary. In England, gatehouses were also growing in size and sophistication until they too challenged the need for a keep in the same castle. The classic Edwardian gatehouse, with two large, flanking towers and multiple portcullises, designed to be defended from attacks both within and outside the main castle, has been often compared to the earlier Norman keeps: some of the largest gatehouses are called gatehouse keeps for this reason.
The quadrangular castle design that emerged in France during the 13th century was another development that removed the need for a keep. Castles had needed additional living space since their first emergence in the 9th century; initially this had been provided by halls in the bailey, then later by ranges of chambers alongside the inside of a bailey wall, such as at Goodrich. But French designs in the late 12th century took the layout of a contemporary unfortified manor house, whose rooms faced around a central, rectangular courtyard, and built a wall around them to form a castle. The result, illustrated initially at Yonne, and later at Château de Farcheville, was a characteristic quadrangular layout with four large, circular corner towers. It lacked a keep, which was not needed to support this design.
Late medieval keeps (14th–16th centuries)
The end of the medieval period saw a fresh resurgence in the building of keeps in western castles. Some castles continued to be built without keeps: the Bastille in the 1370s, for example, combined a now traditional quadrangular design with machicolated corner towers, gatehouses and moat; the walls, innovatively, were of equal height to the towers. This fashion became copied across French and in England, particularly amongst the nouveau riche, for example at Nunney. The royalty and the very wealthiest in France, England and Spain, however, began to construct a small number of keeps on a much larger scale than before, in England sometimes termed tower keeps, as part of new palace fortresses. This shift reflected political and social pressures, such as the desire of the wealthiest lords to have privacy from their growing households of retainers, as well as the various architectural ideas being exchanged across the region, despite the ongoing Hundred Years War between France and England.
The resurgence in French keep design began after the defeat of the royal armies at the battles of Crécy in 1346 and Poitiers in 1356, which caused high levels of social unrest across the remaining French territories. Charles V of France attempted to restore French royal authority and prestige through the construction of a new range of castles. The Château de Vincennes, where a new keep was completed under Charles by 1380, was the first example of these palace fortresses. The keep at Vincennes was highly innovative: six stories high, with a chemin de ronde running around the machicolated battlements; the luxuriously appointed building was protected by an enceinte wall that formed a "fortified envelope" around the keep. The Vincennes keep was copied elsewhere across France, particularly as the French kings reconquered territories from the English, encouraging a style that emphasised very tall keeps with prominent machicolations. No allowance for the emerging new gunpowder weapons was made in these keeps, although later in the century gunports were slowly being added, as for example by Charles VI to his keep at Saint-Malo.
The French model spread into Iberia in the second half of the century, where the most powerful nobles in Castile built a number of similar tall keeps, such as that at Peñafiel, taking advantage of the weakness of the Castilian Crown during the period. Henry IV of Castile responded in the 15th century by creating a sequence of royal castles with prominent keeps at the Castle of La Mota, Portillo, and Alcázar of Segovia: built to particular proportions, these keeps became known as a key element of the Valladolid school of Spanish castle design. Smaller versions of these keeps were subsequently built by many aspiring new aristocracy in Spain, including many converted Jews, keen to improve their social prestige and position in society. The French model of tall keeps was also echoed in some German castles, such as that at Karlštejn, although the layout and positioning of these towers still followed the existing bergfried model, rather than that in western castles. An other impressive 15th century metiterenian castle keep is the keep of the Kolossi Castle, in Cyprus, a three floor square keep, 21 meters high.
The 15th and 16th centuries saw a small number of English and occasional Welsh castles develop still grander keeps. The first of these large tower keeps were built in the north of England during the 14th century, at locations such as Warkworth. They were probably partially inspired by designs in France, but they also reflected the improvements in the security along the Scottish border during the period, and the regional rise of major noble families such as the Percies and the Nevilles, whose wealth encouraged a surge in castle building at the end of the 14th century. New castles at Raby, Bolton, and Warkworth Castle took the quadrangular castle styles of the south and combined them with exceptionally large tower keeps to form a distinctive, northern style. Built by major noble houses, these castles were typically even more opulent than the smaller castles like Nunney, built by the nouveau riche. They marked what historian Anthony Emery has described as a "...second peak of castle building in England and Wales," following on from the Edwardian designs at the end of the 14th century.
In the 15th century, the fashion for the creation of very expensive, French-influenced palatial castles featuring complex tower keeps spread, with new keeps being built at Wardour, Tattershall, and Raglan Castle. In central and eastern England, some keeps began to be built in brick, with Caister and Tattershall forming examples of this trend. In Scotland, the construction of Holyrood Great Tower between 1528 and 1532 drew on this English tradition, but incorporated additional French influences to produce a highly secure but comfortable keep, guarded by a gun park. These tower keeps were expensive buildings to construct, each built to a unique design for a specific lord and, as historian Norman Pounds has suggested, they "...were designed to allow very rich men to live in luxury and splendour."
At the same time as these keeps were being built by the extremely wealthy, much smaller, keep-like structures called tower houses or peel towers were built across Ireland, Scotland, and northern England, often by relatively poorer local lords and landowners. It was originally argued that Irish tower houses were based on the Scottish design, but the pattern of development of such castles in Ireland does not support this hypothesis. A tower house would typically be a tall, square, stone-built, crenelated building; Scottish and Ulster tower houses were often also surrounded by a barmkyn or bawn wall. Most academics have concluded that tower houses should not be classified as keeps but rather as a form of fortified house.
As the 16th century progressed, keeps fell out of fashion once again. In England, the gatehouse also began to supplant the keep as the key focus for a new castle development. By the 15th century, it was increasingly unusual for a lord to build both a keep and a large gatehouse at the same castle, and by the early 16th century, the gatehouse had easily overtaken the keep as the more fashionable feature: indeed, almost no new keeps were built in England after this period. The classical Palladian style began to dominate European architecture during the 17th century, causing a further move away from the use of keeps. Buildings in this style usually required considerable space for the enfiladed formal rooms that became essential for modern palaces by the middle of the century, and this style was impossible to fit into a traditional keep. The keep at Bolsover Castle in England was one of the few to be built as part of a Palladian design.
Later use and destruction of keeps (17th–21st centuries)
From the 17th century onwards, some keeps were deliberately destroyed. In England, many were destroyed after the end of the Second English Civil War in 1649, when Parliament took steps to prevent another royalist uprising by slighting, or damaging, castles so as to prevent them from having any further military utility. Slighting was quite expensive and took considerable effort to carry out, so damage was usually done in the most cost-efficient fashion with only selected walls being destroyed. Keeps were singled out for particular attention in this process because of their continuing political and cultural importance, and the prestige they lent their former royalist owners – at Kenilworth, for example, only the keep was slighted, and at Raglan, the keep was the main focus of parliamentary activity. There was some equivalent destruction of keeps in France in the 17th and 18th centuries, such as the slighting of Montaiguillon by Cardinal Richelieu in 1624, but the catalogue of damage was far less than that of the 1640s and early 1650s in England.
In England, ruined medieval castles became fashionable again in the middle of the 18th century. They were considered an interesting counterpoint to Palladian classical architecture, and gave a degree of medieval allure to their owners. Some keeps were modified to exaggerate this effect: Hawarden, for example, was remodelled to appear taller but also more decayed, the better to produce a good silhouette. The interest continued and, in the late 18th and 19th century, it became fashionable to build intact, replica castles in England, resulting in what A. Rowan has called the Norman style of new castle building, characterised by the inclusion of large keeps; the final replica keep to be built in this way was at Penrhyn between 1820 and 1840.
Where there was an existing castle on a site, another response across 19th-century Europe was to attempt to improve the buildings, bringing their often chaotic historic features into line with a more integrated architectural aesthetic, in a style often termed Gothic Revivalism. There were numerous attempts to restore or rebuild keeps so as to produce this consistently Gothic style: in England, the architect Anthony Salvin was particularly prominent – as illustrated by reworking and heightening of the keep at Windsor Castle, while in France, Eugène Viollet-le-Duc reworked the keeps at castles in locations like Pierrefonds during the 1860s and 1870s, admittedly in a largely speculative fashion, since the original keep had been mostly destroyed in 1617.
The Spanish Civil War and First and Second World Wars in the 20th century caused damage to many castle keeps across Europe; in particular, the famous keep at Coucy was destroyed by the German Army in 1917. By the late 20th century the conservation of castle keeps formed part of government policy across France, England, Ireland, and Spain. In the 21st century in England, most keeps are in ruins and form part of the tourism and heritage industries, rather than being used as functioning buildings – the keep of Windsor Castle being a rare exception. In Germany, large numbers of the bergfried towers were restored as functional buildings in the late 19th and early 20th centuries, often as government offices or youth hostels, or the modern conversion of tower houses, which in many cases have become modernised domestic homes.
| Technology | Fortification | null |
957663 | https://en.wikipedia.org/wiki/Cyclobutane | Cyclobutane | Cyclobutane is a cycloalkane and organic compound with the formula (CH2)4. Cyclobutane is a colourless gas and is commercially available as a liquefied gas. Derivatives of cyclobutane are called cyclobutanes. Cyclobutane itself is of no commercial or biological significance, but more complex derivatives are important in biology and biotechnology.
Structure
The bond angles between carbon atoms are significantly strained and as such have lower bond energies than related linear or unstrained hydrocarbons, e.g. butane or cyclohexane. As such, cyclobutane is unstable above about 500 °C.
The four carbon atoms in cyclobutane are not coplanar; instead, the ring typically adopts a folded or "puckered" conformation. This implies that the C-C-C angle is less than 90°. One of the carbon atoms makes a 25° angle with the plane formed by the other three carbons. In this way, some of the eclipsing interactions are reduced. The conformation is also known as a "butterfly". Equivalent puckered conformations interconvert:
Cyclobutanes in biology and biotechnology
Despite the inherent strain, the cyclobutane motif is indeed found in nature. One example is pentacycloanammoxic acid, which is a ladderane composed of 5 fused cyclobutane units. The estimated strain in this compound is 3 times that of cyclobutane. The compound is found in bacteria performing the anammox process where it forms part of a tight and very dense membrane believed to protect the organism from toxic hydroxylamine and hydrazine involved in the production of nitrogen and water from nitrite ions and ammonia. Some related fenestranes are also found in nature.
Cyclobutane photo dimers (CPD) are formed by photochemical reactions that result in the coupling of the C=C double bonds of pyrimidines. Thymine dimers (T-T dimers) formed in between two thymines are the most abundant of the CPDs. CPDs are readily repaired by nucleotide excision repair enzymes. In most organisms, they can also be repaired by photolyases, a light-dependent family of enzymes. Xeroderma pigmentosum is a genetic disease where this damage can not be repaired, resulting in skin discolouration and tumours induced by exposure to UV light.
Carboplatin is a popular anticancer drug that is derived from cyclobutane-1,1-dicarboxylic acid.
Preparation
Many methods exist for the preparation of cyclobutanes. Alkenes dimerize upon irradiation with UV-light. 1,4-Dihalobutanes convert to cyclobutanes upon dehalogenation with reducing metals.
Cyclobutane was first synthesized in 1907 by James Bruce and Richard Willstätter by hydrogenating cyclobutene in the presence of nickel.
| Physical sciences | Aliphatic hydrocarbons | Chemistry |
30696992 | https://en.wikipedia.org/wiki/Cardboard | Cardboard | Cardboard is a generic term for heavy paper-based products. The construction can range from a thick paper known as paperboard to corrugated fiberboard which is made of multiple plies of material. Natural cardboards can range from grey to light brown in color, depending on the specific product; dyes, pigments, printing, and coatings are available.
The term "cardboard" has general use in English and French, but the term cardboard is deprecated in commerce and industry as not adequately defining a specific product. Material producers, container manufacturers, packaging engineers,
and standards organizations, use more specific terminology.
Usage statistics
In 2020, the United States hit a record high in its yearly use of one of the most ubiquitous manufactured materials on earth, cardboard. With around 80 percent of all the products sold in the United States being packaged in cardboard, over 120 billion pieces were used that year. In the same year, over 13,000 separate pieces of consumer cardboard packaging were thrown away by American households, combined with all paper products, and this constitutes almost 42 percent of all solid waste generated by the United States annually. In an effort to reduce this environmental impact, many households have started repurposing cardboard boxes for eco-friendly purposes.
However, despite the sheer magnitude of paper waste, the vast majority of it is composed of one of the most successful and sustainable packaging materials of modern times - corrugated cardboard, known industrially as corrugated fiberboard.
Types
Various card stocks
Various types of cards are available, which may be called "cardboard". Included are: thick paper (of various types) or pasteboard used for business cards, aperture cards, postcards, playing cards, catalog covers, binder's board for bookbinding, scrapbooking, and other uses which require higher durability than regular paper.
Paperboard
Paperboard is a paper-based material, usually more than about ten mils () thick. It is often used for folding cartons, set-up boxes, carded packaging, etc. Configurations of paperboard include:
Containerboard, used in the production of corrugated fiberboard.
Folding boxboard, comprising multiple layers of chemical and mechanical pulp.
Solid bleached board, made purely from bleached chemical pulp and usually has a mineral or synthetic pigment.
Solid unbleached board, typically made of unbleached chemical pulp.
White lined chipboard, typically made from layers of waste paper or recycled fibers, most often with two to three layers of coating on the top and one layer on the reverse side. Because of its recycled content it will be grey from the inside.
Binder's board, a paperboard used in bookbinding for making hardcovers.
Currently, materials falling under these names may be made without using any actual paper.
Corrugated fiberboard
Corrugated fiberboard is a combination of paperboards, usually two flat liners and one inner fluted corrugated medium. It is often used for making corrugated boxes for shipping or storing products. This type of cardboard is also used by artists as original material for sculpting.
Recycling
Most types of cardboard are recyclable. Boards that are laminates, wax coated, or treated for wet-strength are often more difficult to recycle. Clean cardboard (i.e., cardboard that has not been subject to chemical coatings) "is usually worth recovering, although often the difference between the value it realizes and the cost of recovery is marginal". Cardboard can be recycled for industrial or domestic use. For example, cardboard may be composted or shredded for animal bedding.
History
The material had been first made in France, in 1751, by a pupil of Réaumur, and was used to reinforce playing cards.
The term cardboard has been used since at least 1848, when Anne Brontë mentioned it in her novel The Tenant of Wildfell Hall. The Kellogg brothers first used paperboard cartons to hold their flaked corn cereal, and later, when they began marketing it to the general public, a heat-sealed bag of wax paper was wrapped around the outside of the box and printed with their brand name. This development marked the origin of the cereal box, though in modern times the sealed bag is plastic and is kept inside the box. The Kieckhefer Container Company, run by John W. Kieckhefer, was another early American packaging industry pioneer. It excelled in the use of fiber shipping containers, particularly the paper milk carton.
Examples of different end use
| Technology | Materials | null |
29171632 | https://en.wikipedia.org/wiki/Marine%20habitat | Marine habitat | A marine habitat is a habitat that supports marine life. Marine life depends in some way on the saltwater that is in the sea (the term marine comes from the Latin mare, meaning sea or ocean). A habitat is an ecological or environmental area inhabited by one or more living species. The marine environment supports many kinds of these habitats.
Marine habitats can be divided into coastal and open ocean habitats. Coastal habitats are found in the area that extends from as far as the tide comes in on the shoreline out to the edge of the continental shelf. Most marine life is found in coastal habitats, even though the shelf area occupies only seven percent of the total ocean area. Open ocean habitats are found in the deep ocean beyond the edge of the continental shelf.
Alternatively, marine habitats can be divided into pelagic and demersal zones. Pelagic habitats are found near the surface or in the open water column, away from the bottom of the ocean. Demersal habitats are near or on the bottom of the ocean. An organism living in a pelagic habitat is said to be a pelagic organism, as in pelagic fish. Similarly, an organism living in a demersal habitat is said to be a demersal organism, as in demersal fish. Pelagic habitats are intrinsically shifting and ephemeral, depending on what ocean currents are doing.
Marine habitats can be modified by their inhabitants. Some marine organisms, like corals, kelp, mangroves and seagrasses, are ecosystem engineers which reshape the marine environment to the point where they create further habitat for other organisms. By volume the ocean provides most of the habitable space on the planet.
Overview
In contrast to terrestrial habitats, marine habitats are shifting and ephemeral. Swimming organisms find areas by the edge of a continental shelf a good habitat, but only while upwellings bring nutrient rich water to the surface. Shellfish find habitat on sandy beaches, but storms, tides and currents mean their habitat continually reinvents itself.
The presence of seawater is common to all marine habitats. Beyond that many other things determine whether a marine area makes a good habitat and the type of habitat it makes. For example:
temperature – is affected by geographical latitude, ocean currents, weather, the discharge of rivers, and by the presence of hydrothermal vents or cold seeps
sunlight – photosynthetic processes depend on how deep and turbid the water is
nutrients – are transported by ocean currents to different marine habitats from land runoff, or by upwellings from the deep sea, or they sink through the sea as marine snow
salinity – varies, particularly in estuaries or near river deltas, or by hydrothermal vents
dissolved gases – oxygen levels in particular, can be increased by wave actions and decreased during algal blooms
acidity – this is partly to do with dissolved gases above, since the acidity of the ocean is largely controlled by how much carbon dioxide is in the water.
turbulence – ocean waves, fast currents and the agitation of water affect the nature of habitats
cover – the availability of cover such as the adjacency of the sea bottom, or the presence of floating objects
substrate – The slope, orientation, profile and rugosity of hard substrates, and particle size, sorting and density of unconsolidated sediment bottoms can make a big difference to the life forms that can settle on it.
the occupying organisms themselves – since organisms modify their habitats by the act of occupying them, and some, like corals, kelp, mangroves and seagrasses, create further habitats for other organisms.
There are five major oceans, of which the Pacific Ocean is nearly as large as the rest put together. Coastlines fringe the land for nearly 380,000 kilometres.
Altogether, the ocean occupies 71 percent of the world surface, averaging nearly four kilometres in depth. By volume, the ocean contains more than 99 percent of the Earth's liquid water. The science fiction writer Arthur C. Clarke has pointed out it would be more appropriate to refer to the planet Earth as the planet Sea or the planet
Ocean.
Marine habitats can be broadly divided into pelagic and demersal habitats. Pelagic habitats are the habitats of the open water column, away from the bottom of the ocean. Demersal habitats are the habitats that are near or on the bottom of the ocean. An organism living in a pelagic habitat is said to be a pelagic organism, as in pelagic fish. Similarly, an organism living in a demersal habitat is said to be a demersal organism, as in demersal fish. Pelagic habitats are intrinsically ephemeral, depending on what ocean currents are doing.
The land-based ecosystem depends on topsoil and fresh water, while the marine ecosystem depends on dissolved nutrients washed down from the land.
Ocean deoxygenation poses a threat to marine habitats, due to the growth of low oxygen zones.
Ocean currents
In marine systems, ocean currents have a key role determining which areas are effective as habitats, since ocean currents transport the basic nutrients needed to support marine life. Plankton are the life forms that inhabit the ocean that are so small (less than 2 mm) that they cannot effectively propel themselves through the water, but must drift instead with the currents. If the current carries the right nutrients, and if it also flows at a suitably shallow depth where there is plenty of sunlight, then such a current itself can become a suitable habitat for photosynthesizing tiny algae called phytoplankton. These tiny plants are the primary producers in the ocean, at the start of the food chain. In turn, as the population of drifting phytoplankton grows, the water becomes a suitable habitat for zooplankton, which feed on the phytoplankton. While phytoplankton are tiny drifting plants, zooplankton are tiny drifting animals, such as the larvae of fish and marine invertebrates. If sufficient zooplankton establish themselves, the current becomes a candidate habitat for the forage fish that feed on them. And then if sufficient forage fish move to the area, it becomes a candidate habitat for larger predatory fish and other marine animals that feed on the forage fish. In this dynamic way, the current itself can, over time, become a moving habitat for multiple types of marine life.
Ocean currents can be generated by differences in the density of the water. How dense water is depends on how saline or warm it is. If water contains differences in salt content or temperature, then the different densities will initiate a current. Water that is saltier or cooler will be denser, and will sink in relation to the surrounding water. Conversely, warmer and less salty water will float to the surface. Atmospheric winds and pressure differences also produces surface currents, waves and seiches. Ocean currents are also generated by the gravitational pull of the sun and moon (tides), and seismic activity (tsunami).
The rotation of the Earth affects the direction ocean currents take, and explains which way the large circular ocean gyres rotate in the image above left. Suppose a current at the equator is heading north. The Earth rotates eastward, so the water possesses that rotational momentum. But the further the water moves north, the slower the earth moves eastward. If the current could get to the North Pole, the earth would not be moving eastward at all. To conserve its rotational momentum, the further the current travels north the faster it must move eastward. So the effect is that the current curves to the right. This is the Coriolis effect. It is weakest at the equator and strongest at the poles. The effect is opposite south of the equator, where currents curve left.
Topography
Biomass
One measure of the relative importance of different marine habitats is the rate at which they produce biomass.
Coastal
Marine coasts are dynamic environments which constantly change, like the ocean which partially shape them. The Earth's natural processes, including weather and sea level change, result in the erosion, accretion and resculpturing of coasts as well as the flooding and creation of continental shelves and drowned river valleys.
The main agents responsible for deposition and erosion along coastlines are waves, tides and currents. The formation of coasts also depends on the nature of the rocks they are made of – the harder the rocks the less likely they are to erode, so variations in rock hardness result in coastlines with different shapes.
Tides often determine the range over which sediment is deposited or eroded. Areas with high tidal ranges allow waves to reach farther up the shore, and areas with lower tidal ranges produce deposition at a smaller elevation interval. The tidal range is influenced by the size and shape of the coastline. Tides do not typically cause erosion by themselves; however, tidal bores can erode as the waves surge up river estuaries from the ocean.
Waves erode coastline as they break on shore releasing their energy; the larger the wave the more energy it releases and the more sediment it moves. Sediment deposited by waves comes from eroded cliff faces and is moved along the coastline by the waves. Sediment deposited by rivers is the dominant influence on the amount of sediment located on a coastline.
The sedimentologist Francis Shepard classified coasts as primary or secondary.
Primary coasts are shaped by non-marine processes, by changes in the land form. If a coast is in much the same condition as it was when sea level was stabilised after the last ice age, it is called a primary coast. "Primary coasts are created by erosion (the wearing away of soil or rock), deposition (the buildup of sediment or sand) or tectonic activity (changes in the structure of the rock and soil because of earthquakes). Many of these coastlines were formed as the sea level rose during the last 18,000 years, submerging river and glacial valleys to form bays and fjords." An example of a primary coast is a river delta, which forms when a river deposits soil and other material as it enters the sea.
Secondary coasts are produced by marine processes, such as the action of the sea or by creatures that live in it. Secondary coastlines include sea cliffs, barrier islands, mud flats, coral reefs, mangrove swamps and salt marshes.
Continental coastlines usually have a continental shelf, a shelf of relatively shallow water, less than 200 metres deep, which extends 68 km on average beyond the coast. Worldwide, continental shelves occupy a total area of about 24 million km2 (9 million sq mi), 8% of the ocean's total area and nearly 5% of the world's total area. Since the continental shelf is usually less than 200 metres deep, it follows that coastal habitats are generally photic, situated in the sunlit epipelagic zone. This means the conditions for photosynthetic processes so important for primary production, are available to coastal marine habitats. Because land is nearby, there are large discharges of nutrient rich land runoff into coastal waters. Further, periodic upwellings from the deep ocean can provide cool and nutrient rich currents along the edge of the continental shelf.
As a result, coastal marine life is the most abundant in the world. It is found in tidal pools, fjords and estuaries, near sandy shores and rocky coastlines, around coral reefs and on or above the continental shelf. Coastal fish include small forage fish as well as the larger predator fish that feed on them. Forage fish thrive in inshore waters where high productivity results from upwelling and shoreline run off of nutrients. Some are partial residents that spawn in streams, estuaries and bays, but most complete their life cycle in the zone. There can also be a mutualism between species that occupy adjacent marine habitats. For example, fringing reefs just below low tide level have a mutually beneficial relationship with mangrove forests at high tide level and sea grass meadows in between: the reefs protect the mangroves and seagrass from strong currents and waves that would damage them or erode the sediments in which they are rooted, while the mangroves and seagrass protect the coral from large influxes of silt, fresh water and pollutants. This additional level of variety in the environment is beneficial to many types of coral reef animals, which for example may feed in the sea grass and use the reefs for protection or breeding.
Coastal habitats are the most visible marine habitats, but they are not the only important marine habitats. Coastlines run for 380,000 kilometres, and the total volume of the ocean is 1,370 million cu km. This means that for each metre of coast, there is 3.6 cu km of ocean space available somewhere for marine habitats.
Intertidal
Intertidal zones, those areas close to shore, are constantly being exposed and covered by the ocean's tides. A huge array of life lives within this zone.
Shore habitats range from the upper intertidal zones to the area where land vegetation takes prominence. It can be underwater anywhere from daily to very infrequently. Many species here are scavengers, living off of sea life that is washed up on the shore. Many land animals also make much use of the shore and intertidal habitats. A subgroup of organisms in this habitat bores and grinds exposed rock through the process of bioerosion.
Sandy shores
Sandy shores, also called beaches, are coastal shorelines where sand accumulates. Waves and currents shift the sand, continually building and eroding the shoreline. Longshore currents flow parallel to the beaches, making waves break obliquely on the sand. These currents transport large amounts of sand along coasts, forming spits, barrier islands and tombolos. Longshore currents also commonly create offshore bars, which give beaches some stability by reducing erosion.
Sandy shores are full of life. The grains of sand host diatoms, bacteria and other microscopic creatures. Some fish and turtles return to certain beaches and spawn eggs in the sand. Birds habitat beaches, like gulls, loons, sandpipers, terns and pelicans. Aquatic mammals, such sea lions, recuperate on them. Clams, periwinkles, crabs, shrimp, starfish and sea urchins are found on most beaches.
Sand is a sediment made from small grains or particles with diameters between about 60 μm and 2 mm. Mud (see mudflats below) is a sediment made from particles finer than sand. This small particle size means that mud particles tend to stick together, whereas sand particles do not. Mud is not easily shifted by waves and currents, and when it dries out, cakes into a solid. By contrast, sand is easily shifted by waves and currents, and when sand dries out it can be blown in the wind, accumulating into shifting sand dunes. Beyond the high tide mark, if the beach is low-lying, the wind can form rolling hills of sand dunes. Small dunes shift and reshape under the influence of the wind while larger dunes stabilise the sand with vegetation.
Ocean processes grade loose sediments to particle sizes other than sand, such as gravel or cobbles. Waves breaking on a beach can leave a berm, which is a raised ridge of coarser pebbles or sand, at the high tide mark. Shingle beaches are made of particles larger than sand, such as cobbles, or small stones. These beaches make poor habitats. Little life survives because the stones are churned and pounded together by waves and currents.
Rocky shores
The relative solidity of rocky shores seems to give them a permanence compared to the shifting nature of sandy shores. This apparent stability is not real over even quite short geological time scales, but it is real enough over the short life of an organism. In contrast to sandy shores, plants and animals can anchor themselves to the rocks.
Competition can develop for the rocky spaces. For example, barnacles can compete successfully on open intertidal rock faces to the point where the rock surface is covered with them. Barnacles resist desiccation and grip well to exposed rock faces. However, in the crevices of the same rocks, the inhabitants are different. Here mussels can be the successful species, secured to the rock with their byssal threads.
Rocky and sandy coasts are vulnerable because humans find them attractive and want to live near them. An increasing proportion of the humans live by the coast, putting pressure on coastal habitats.
Mudflats
Mudflats are coastal wetlands that form when mud is deposited by tides or rivers. They are found in sheltered areas such as bays, bayous, lagoons, and estuaries. Mudflats may be viewed geologically as exposed layers of bay mud, resulting from deposition of estuarine silts, clays and marine animal detritus. Most of the sediment within a mudflat is within the intertidal zone, and thus the flat is submerged and exposed approximately twice daily.
Mangrove forests and salt marshes
Mangrove swamps and salt marshes form important coastal habitats in tropical and temperate areas respectively.
Mangroves are species of shrubs and medium size trees that grow in saline coastal sediment habitats in the tropics and subtropics – mainly between latitudes ° N and ° S. The saline conditions tolerated by various species range from brackish water, through pure seawater (30 to 40 ppt), to water concentrated by evaporation to over twice the salinity of ocean seawater (up to 90 ppt). There are many mangrove species, not all closely related. The term "mangrove" is used generally to cover all of these species, and it can be used narrowly to cover just mangrove trees of the genus Rhizophora.
Mangroves form a distinct characteristic saline woodland or shrubland habitat, called a mangrove swamp or mangrove forest. Mangrove swamps are found in depositional coastal environments, where fine sediments (often with high organic content) collect in areas protected from high-energy wave action. Mangroves dominate three quarters of tropical coastlines.
Estuaries
An estuary is a partly enclosed coastal body of water with one or more rivers or streams flowing into it, and with a free connection to the open sea. Estuaries form a transition zone between river environments and ocean environments and are subject to both marine influences, such as tides, waves, and the influx of saline water; and riverine influences, such as flows of fresh water and sediment. The inflow of both seawater and freshwater provide high levels of nutrients in both the water column and sediment, making estuaries among the most productive natural habitats in the world.
Most estuaries were formed by the flooding of river-eroded or glacially scoured valleys when sea level began to rise about 10,000-12,000 years ago. They are amongst the most heavily populated areas throughout the world, with about 60% of the world's population living along estuaries and the coast. As a result, estuaries are suffering degradation by many factors, including sedimentation from soil erosion from deforestation; overgrazing and other poor farming practices; overfishing; drainage and filling of wetlands; eutrophication due to excessive nutrients from sewage and animal wastes; pollutants including heavy metals, PCBs, radionuclides and hydrocarbons from sewage inputs; and diking or damming for flood control or water diversion.
Estuaries provide habitats for a large number of organisms and support very high productivity. Estuaries provide habitats for salmon and sea trout nurseries, as well as migratory bird populations. Two of the main characteristics of estuarine life are the variability in salinity and sedimentation. Many species of fish and invertebrates have various methods to control or conform to the shifts in salt concentrations and are termed osmoconformers and osmoregulators. Many animals also burrow to avoid predation and to live in the more stable sedimental environment. However, large numbers of bacteria are found within the sediment which have a very high oxygen demand. This reduces the levels of oxygen within the sediment often resulting in partially anoxic conditions, which can be further exacerbated by limited water flux. Phytoplankton are key primary producers in estuaries. They move with the water bodies and can be flushed in and out with the tides. Their productivity is largely dependent on the turbidity of the water. The main phytoplankton present are diatoms and dinoflagellates which are abundant in the sediment.
Kelp forests
Kelp forests are underwater areas with a high density of kelp. They form some of the most productive and dynamic ecosystems on Earth. Smaller areas of anchored kelp are called kelp beds. Kelp forests occur worldwide throughout temperate and polar coastal oceans.
Kelp forests provide a unique three-dimensional habitat for marine organisms and are a source for understanding many ecological processes. Over the last century, they have been the focus of extensive research, particularly in trophic ecology, and continue to provoke important ideas that are relevant beyond this unique ecosystem. For example, kelp forests can influence coastal oceanographic patterns and provide many ecosystem services.
However, humans have contributed to kelp forest degradation. Of particular concern are the effects of overfishing nearshore ecosystems, which can release herbivores from their normal population regulation and result in the over-grazing of kelp and other algae. This can rapidly result in transitions to barren landscapes where relatively few species persist.
Frequently considered an ecosystem engineer, kelp provides a physical substrate and habitat for kelp forest communities. In algae (Kingdom: Protista), the body of an individual organism is known as a thallus rather than as a plant (Kingdom: Plantae). The morphological structure of a kelp thallus is defined by three basic structural units:
The holdfast is a root-like mass that anchors the thallus to the sea floor, though unlike true roots it is not responsible for absorbing and delivering nutrients to the rest of the thallus;
The stipe is analogous to a plant stalk, extending vertically from the holdfast and providing a support framework for other morphological features;
The fronds are leaf- or blade-like attachments extending from the stipe, sometimes along its full length, and are the sites of nutrient uptake and photosynthetic activity.
In addition, many kelp species have pneumatocysts, or gas-filled bladders, usually located at the base of fronds near the stipe. These structures provide the necessary buoyancy for kelp to maintain an upright position in the water column.
The environmental factors necessary for kelp to survive include hard substrate (usually rock), high nutrients (e.g., nitrogen, phosphorus), and light (minimum annual irradiance dose > 50 E m−2). Especially productive kelp forests tend to be associated with areas of significant oceanographic upwelling, a process that delivers cool nutrient-rich water from depth to the ocean's mixed surface layer. Water flow and turbulence facilitate nutrient assimilation across kelp fronds throughout the water column. Water clarity affects the depth to which sufficient light can be transmitted. In ideal conditions, giant kelp (Macrocystis spp.) can grow as much as 30-60 centimetres vertically per day. Some species such as Nereocystis are annual while others like Eisenia are perennial, living for more than 20 years. In perennial kelp forests, maximum growth rates occur during upwelling months (typically spring and summer) and die-backs correspond to reduced nutrient availability, shorter photoperiods and increased storm frequency.
Seagrass meadows
Seagrasses are flowering plants from one of four plant families which grow in marine environments. They are called seagrasses because the leaves are long and narrow and are very often green, and because the plants often grow in large meadows which look like grassland. Since seagrasses photosynthesize and are submerged, they must grow submerged in the photic zone, where there is enough sunlight. For this reason, most occur in shallow and sheltered coastal waters anchored in sand or mud bottoms.
Seagrasses form extensive beds or meadows, which can be either monospecific (made up of one species) or multispecific (where more than one species co-exist). Seagrass beds make highly diverse and productive ecosystems. They are home to phyla such as juvenile and adult fish, epiphytic and free-living macroalgae and microalgae, mollusks, bristle worms, and nematodes. Few species were originally considered to feed directly on seagrass leaves (partly because of their low nutritional content), but scientific reviews and improved working methods have shown that seagrass herbivory is a highly important link in the food chain, with hundreds of species feeding on seagrasses worldwide, including green turtles, dugongs, manatees, fish, geese, swans, sea urchins and crabs.
Seagrasses are ecosystem engineers in the sense that they partly create their own habitat. The leaves slow down water-currents increasing sedimentation, and the seagrass roots and rhizomes stabilize the seabed. Their importance to associated species is mainly due to provision of shelter (through their three-dimensional structure in the water column), and due to their extraordinarily high rate of primary production. As a result, seagrasses provide coastal zones with ecosystem services, such as fishing grounds, wave protection, oxygen production and protection against coastal erosion. Seagrass meadows account for 15% of the ocean's total carbon storage.
Reefs
A reef is a ridge or shoal of rock, coral or similar relatively stable material, lying beneath the surface of a natural body of water. Many reefs result from natural, abiotic processes but there are also reefs such as the coral reefs of tropical waters formed by biotic processes dominated by corals and coralline algae. Artificial reefs such as shipwrecks and other anthropogenic underwater structures may occur intentionally or as the result of an accident, and sometimes have a designed role in enhancing the physical complexity of featureless sand bottoms, thereby attracting a more diverse assemblage of organisms. Reefs are often quite near to the surface, but not all definitions require this. Fringing reefs, the most common type of reef, are found close to shorelines and surrounding islands.
Rocky reefs
Coral reefs
Coral reefs comprise some of the densest and most diverse habitats in the world. The best-known types of reefs are tropical coral reefs which exist in most tropical waters; however, coral reefs can also exist in cold water. Reefs are built up by corals and other calcium-depositing animals, usually on top of a rocky outcrop on the ocean floor. Reefs can also grow on other surfaces, which has made it possible to create artificial reefs. Coral reefs also support a huge community of life, including the corals themselves, their symbiotic zooxanthellae, tropical fish and many other organisms.
Much attention in marine biology is focused on coral reefs and the El Niño weather phenomenon. In 1998, coral reefs experienced the most severe mass bleaching events on record, when vast expanses of reefs across the world died because sea surface temperatures rose well above normal. Some reefs are recovering, but scientists say that between 50% and 70% of the world's coral reefs are now endangered and predict that global warming could exacerbate this trend.
Surface waters
Surface microlayer
The surface microlayer of the ocean serves as the transitional area between the atmosphere and the ocean. It covers around 70% of the Earth's surface as it covers most of the ocean waters on the planet. The microlayer is known for its unique biological and chemical properties which give it a small ecosystem of its own and serves as a distinct habitat from the deeper ocean waters.
The surface microlayer is not in fact entirely aqueous like the rest of the ocean, but is closer to a kind of hydrated gel composed of concentrated nutrients forming a biological film over the water it covers. This film is rich in microbes which mediate the interactions between the sun, the atmosphere, and the waters below.
Although thin, the surface microlayer is critical for life beneath it. Because of the environment rich in microbes and nutrients, larvae of fish and other aquatic animals are often laid in the microlayer to incubate. The plankton in the microlayer are distinctly adapted to withstand high levels of radiation, and serve as buffers to prevent this potentially harmful radiation from reaching the deeper water. Environmental changes such as aerosols or dust storms can cause these surface plankton to become overproductive, leading to blooms.
Because of the unique properties of the microlayer, pollutants often accumulate within and use it to reach other parts of the ocean. Hydrophobic compounds, such as petroleum, flame retardants, and heavy metals, have a particular affinity for the surface microlayer. Recently, the abundance of aerosols and microplastics has also had an impact on the SML and their accumulation has led to many problems, such as animal ingestion of these compounds leading to widespread disruption of balance and spread of these compounds among marine communities.
The surface microlayer is also critical to gas exchange between the atmosphere and the ocean. Because the microlayer is filled with microbes, it is widely theorized that it plays a critical role in gas exchange and uptake of nutrients, but relatively little data on this has been collected. The central feature of the microlayer is the temperature, as it is an indicator of how pollutants and human activity affects the ocean.
Epipelagic zone
The surface waters are sunlit. The waters down to about 200 metres are said to be in the epipelagic zone. Enough sunlight enters the epipelagic zone to allow photosynthesis by phytoplankton. The epipelagic zone is usually low in nutrients. This partially because the organic debris produced in the zone, such as excrement and dead animals, sink to the depths and are lost to the upper zone. Photosynthesis can happen only if both sunlight and nutrients are present.
In some places, like at the edge of continental shelves, nutrients can upwell from the ocean depth, or land runoff can be distributed by storms and ocean currents. In these areas, given that both sunlight and nutrients are now present, phytoplankton can rapidly establish itself, multiplying so fast that the water turns green from the chlorophyll, resulting in an algal bloom. These nutrient rich surface waters are among the most biologically productive in the world, supporting billions of tonnes of biomass.
"Phytoplankton are eaten by zooplankton - small animals which, like phytoplankton, drift in the ocean currents. The most abundant zooplankton species are copepods and krill: tiny crustaceans that are the most numerous animals on Earth. Other types of zooplankton include jelly fish and the larvae of fish, marine worms, starfish, and other marine organisms". In turn, the zooplankton are eaten by filter-feeding animals, including some seabirds, small forage fish like herrings and sardines, whale sharks, manta rays, and the largest animal in the world, the blue whale. Yet again, moving up the foodchain, the small forage fish are in turn eaten by larger predators, such as tuna, marlin, sharks, large squid, seabirds, dolphins, and toothed whales.
Open ocean
The open ocean is relatively unproductive because of a lack of nutrients, yet because it is so vast, it has more overall primary production than any other marine habitat. Only about 10 percent of marine species live in the open ocean. But among them are the largest and fastest of all marine animals, as well as the animals that dive the deepest and migrate the longest. In the depths lurk animal that, to our eyes, appear hugely alien.
Deep sea
The deep sea starts at the aphotic zone, the point where sunlight loses most of its energy in the water. Many life forms that live at these depths have the ability to create their own light a unique evolution known as bio-luminescence.
In the deep ocean, the waters extend far below the epipelagic zone, and support very different types of pelagic life forms adapted to living in these deeper zones.
Much of the aphotic zone's energy is supplied by the open ocean in the form of detritus. In deep water, marine snow is a continuous shower of mostly organic detritus falling from the upper layers of the water column. Its origin lies in activities within the productive photic zone. Marine snow includes dead or dying plankton, protists (diatoms), fecal matter, sand, soot and other inorganic dust. The "snowflakes" grow over time and may reach several centimetres in diameter, travelling for weeks before reaching the ocean floor. However, most organic components of marine snow are consumed by microbes, zooplankton and other filter-feeding animals within the first 1,000 metres of their journey, that is, within the epipelagic zone. In this way marine snow may be considered the foundation of deep-sea mesopelagic and benthic ecosystems: As sunlight cannot reach them, deep-sea organisms rely heavily on marine snow as an energy source.
Some deep-sea pelagic groups, such as the lanternfish, ridgehead, marine hatchetfish, and lightfish families are sometimes termed pseudoceanic because, rather than having an even distribution in open water, they occur in significantly higher abundances around structural oases, notably seamounts and over continental slopes. The phenomenon is explained by the likewise abundance of prey species which are also attracted to the structures.
The fish in the different pelagic and deep water benthic zones are physically structured, and behave in ways, that differ markedly from each other. Groups of coexisting species within each zone all seem to operate in similar ways, such as the small mesopelagic vertically migrating plankton-feeders, the bathypelagic anglerfishes, and the deep water benthic rattails. "
Ray finned species, with spiny fins, are rare among deep sea fishes, which suggests that deep sea fish are ancient and so well adapted to their environment that invasions by more modern fishes have been unsuccessful. The few ray fins that do exist are mainly in the Beryciformes and Lampriformes, which are also ancient forms. Most deep sea pelagic fishes belong to their own orders, suggesting a long evolution in deep sea environments. In contrast, deep water benthic species, are in orders that include many related shallow water fishes.
The umbrella mouth gulper is a deep sea eel with an enormous loosely hinged mouth. It can open its mouth wide enough to swallow a fish much larger than itself, and then expand its stomach to accommodate its catch.
Sea floor
Vents and seeps
Hydrothermal vents along the mid-ocean ridge spreading centers act as oases, as do their opposites, cold seeps. Such places support unique marine biomes and many new marine microorganisms and other lifeforms have been discovered at these locations.
Trenches
The deepest recorded oceanic trenches measure to date is the Mariana Trench, near the Philippines, in the Pacific Ocean at 10,924 m (35,838 ft). At such depths, water pressure is extreme and there is no sunlight, but some life still exists. A white flatfish, a shrimp and a jellyfish were seen by the American crew of the bathyscaphe Trieste when it dove to the bottom in 1960.
Seamounts
Marine life also flourishes around seamounts that rise from the depths, where fish and other sea life congregate to spawn and feed.
Anthropogenic impacts
Mudflats are typically important regions for wildlife, supporting a large population, although levels of biodiversity are not particularly high. They are of particular importance to migratory birds as well as crabs, shrimp, and shellfish. These areas along the coast act as a nursery for these animals by providing an area for reproduction and feeding. However, this can pose as an issue due to the high trafficking of the birds migrating for nesting, then leaving to return to their seasonal homes. Whatever pollutants the birds take in while breeding are brought back with them to their next location, thus polluting that area as well. In the United Kingdom mudflats have been classified as a Biodiversity Action Plan priority habitat. European countries such as France have also found it beneficial to use the Marine Influence Index (MII) to be able to monitor the responses to pollution the local plant and animal species may have as well as monitor any type of deviation from the natural patterns displayed previously.
Although many parts of the seafloor have yet to be explored, researchers have found that parts of it have been greatly affected by human activity. Bottom trawling, microplastic pollution, and industrial metals have slowly changed and altered the composition of the sea floor. Bottom trawling refers to a commercial deep sea fishing technique in which the equipment drags across the sea floor. This has had an adverse effect on the seafloor as it changes the surface structure and composition. In addition, microplastic pollution has become an increasing problem to the seafloor as plastics and other debris are found in many of the sediments. Due to the build up of litter, the habitats and environments of organisms on the seafloor are being impacted and changed. This includes industrial facilities dumping new metals and minerals, such as cadmium, onto the seafloor that change the chemical composition of the water and poison the inhabitants.
There are also negative anthropogenic impacts on deep sea habitats, including trash pollution and chemical pollution. Plastic pollution in particular, is one of the greatest forms of uncontrolled human activity that is visible in our oceans today. Researchers in the Northwestern south China Sea recorded large plastic-dominated litter piles in submarine canyons. These durable plastics can diffuse into smaller organisms and are then inadvertently consumed by humans in the food we eat and water we drink. Another threat to organisms lurking in the deep ocean is ghost fishing, and bycatch. Ghost fishing is the term that refers to any abandoned fishing gear in the ocean that continues to entangle and trap marine organisms. Gill nets for example, have been recorded tangled around deep sea corals and continue ghost fishing for extended periods of time.
Gallery
| Physical sciences | Oceanography | Earth science |
2107349 | https://en.wikipedia.org/wiki/Heterothermy | Heterothermy | Heterothermy or heterothermia (from Greek ἕτερος heteros "other" and θέρμη thermē "heat") is a physiological term for animals that vary between self-regulating their body temperature, and allowing the surrounding environment to affect it. In other words, they exhibit characteristics of both poikilothermy and homeothermy.
Definition
Heterothermic animals are those that can switch between poikilothermic and homeothermic strategies. These changes in strategies typically occur on a daily basis or on an annual basis. More often than not, it is used as a way to dissociate the fluctuating metabolic rates seen in some small mammals and birds (e.g. bats and hummingbirds), from those of traditional cold blooded animals. In many bat species, body temperature and metabolic rate are elevated only during activity. When at rest, these animals reduce their metabolisms drastically, which results in their body temperature dropping to that of the surrounding environment. This makes them homeothermic when active, and poikilothermic when at rest. This phenomenon has been termed 'daily torpor' and was intensively studied in the Djungarian hamster. During the hibernation season, this animal shows strongly reduced metabolism each day during the rest phase while it reverts to endothermic metabolism during its active phase, leading to normal euthermic body temperatures (around 38 °C).
Larger mammals (e.g. ground squirrels) and bats show multi-day torpor bouts during hibernation (up to several weeks) in winter. During these multi-day torpor bouts, body temperature drops to ~1 °C above ambient temperature and metabolism may drop to about 1% of the normal endothermic metabolic rate. Even in these deep hibernators, the long periods of torpor is interrupted by bouts of endothermic metabolism, called arousals (typically lasting between 4–20 hours). These metabolic arousals cause body temperature to return to euthermic levels 35-37 °C. Most of the energy spent during hibernation is spent in arousals (70-80%), but their function remains unresolved.
Shallow hibernation patterns without arousals have been described in large mammals (like the black bear,) or under special environmental circumstances.
Regional heterothermy
Regional heterothermy describes organisms that are able to maintain different temperature "zones" in different regions of the body. This usually occurs in the limbs, and is made possible through the use of counter-current heat exchangers, such as the rete mirabile found in tuna and certain birds. These exchangers equalize the temperature between hot arterial blood going out to the extremities and cold venous blood coming back, thus reducing heat loss. Penguins and many arctic birds use these exchangers to keep their feet at roughly the same temperature as the surrounding ice. This keeps the birds from getting stuck on an ice sheet. Other animals, like the leatherback sea turtle, use the heat exchangers to gather, and retain heat generated by their muscular flippers. There are even some insects which possess this mechanism (see insect thermoregulation), the best-known example being bumblebees, which exhibit counter-current heat exchange at the point of constriction between the mesosoma ("thorax") and metasoma ("abdomen"); heat is retained in the thorax and lost from the abdomen. Using a very similar mechanism, the internal temperature of a honeybee's thorax can exceed 45 °C while in flight.
| Biology and health sciences | Basics | Biology |
2108090 | https://en.wikipedia.org/wiki/Pearlfish | Pearlfish | Pearlfish are marine fish in the ray-finned fish family Carapidae. Pearlfishes inhabit the tropical waters of the Atlantic, Indian, and Pacific Oceans at depths to , along oceanic shelves and slopes. They are slender, elongated fish with no scales, translucent bodies, and dorsal fin rays which are shorter than their anal fin rays. Adults of most species live symbiotically inside various invertebrate hosts, and some live parasitically inside sea cucumbers. The larvae are free living.
Characteristics
Pearlfishes are slender, distinguished by having dorsal fin rays that are shorter than their anal fin rays. They have translucent, scaleless bodies reminiscent of eels. The largest pearlfish are about in length. They reproduce by laying oval-shaped eggs, about 1 mm in length.
Ecology
Pearlfishes are unusual in that the adults of most species live inside various types of invertebrates. They typically live inside clams, sea cucumbers, starfish, or sea squirts, and are simply commensal, not harming their hosts. However, some species are known to be parasitic towards sea cucumbers, eating their gonads and other internal organs. Pearlfish usually live alone, or in pairs.
Regardless of the habits of the adults, the larvae of pearlfish are free-living among the plankton. Pearlfish larvae can be distinguished by the presence of a long filament in front of their dorsal fins, sometimes with various appendages attached.
Genera
The genera are divided into three major groupings based on their level of symbiosis:
Echiodon and Snyderidia - free-living
Carapus and Onuxodon - commensal
Encheliophis - parasitic, fish in this group live in invertebrate hosts found in shallow-water coral communities such as bivalves, sea cucumbers, and starfish.
| Biology and health sciences | Acanthomorpha | Animals |
2109865 | https://en.wikipedia.org/wiki/Dolomedes | Dolomedes | Dolomedes is a genus of large spiders of the family Dolomedidae. They are also known as fishing spiders, raft spiders, dock spiders or wharf spiders. Almost all Dolomedes species are semiaquatic, with the exception of the tree-dwelling D. albineus of the southeastern United States. Many species have a striking pale stripe down each side of the body.
They hunt by waiting at the edge of a pool or stream, then when they detect the ripples from prey, they run across the surface to subdue it using their foremost legs, which are tipped with small claws; then injecting venom with their hollow chelicerae to kill and digest the prey. They mainly eat insects, but some larger species are able to catch small fish. They can also climb beneath the water, when they become encased in a silvery film of air. "Dolomedes" is derived from the Greek word which means wily, deceitful.
There are over a hundred species of Dolomedes throughout the world; examples include Dolomedes aquaticus, a forest-stream species of New Zealand, the raft spider (D. fimbriatus), which lives in bogs in Europe, and the great raft spider (D. plantarius), which lives in fens, also in Europe. Many species are large, some with females up to long with a leg span of .
Aquatic adaptations
Dolomedes spiders are covered all over in short, velvety hairs which are hydrophobic. This allows them to use surface tension to stand or run on the water, like water striders. They can also climb beneath the water, and then air becomes trapped in the body hairs and forms a thin film over the whole surface of the body and legs, giving them the appearance of fine polished silver. Like other spiders, Dolomedes breathe with book lungs beneath their abdomens, and these open into the air film, allowing the spiders to breathe while submerged. The trapped air makes them very buoyant and if they do not hold onto a rock or a plant stem they float to the surface where they pop onto the surface film, completely dry.
Identification
If any of this species are seen without context, one may confuse them with the family Lycosidae, otherwise known as wolf spiders. They can be differentiated from wolf spiders by their smaller posterior median eyes, and eye arrangement of eight eyes in two rows as opposed to three rows. If this is insufficient, one can further differentiate them thanks to their aquatic adaptations.
Hunting behavior
Rather than hunting on land or by waiting in a web, these spiders hunt on the surface of the water itself, preying on mayflies, other aquatic insects, and even small fish. For fishing spiders, the water surface serves the same function as a web does for other spiders. They extend their legs onto the surface, feeling for vibrations given off by prey.
Dolomedes are nocturnal hunters, feeding when birds, their main predators, are sleeping. The method they use to fish for insects is to hold on to the shore with their back legs while the rest of their body lies on the water, with legs stretched out. Dolomedes species tend to be robust with thickset legs that allow them to tackle prey larger than themselves. They stretch out their front legs and wait, as if listening. Their front legs feel the vibrations carried on the water, just as other spiders feel the vibrations in a web. They are able to tell what is causing the vibrations that the water is carrying – to distinguish the drawn-out, erratic vibrations of a struggling insect from the one-off vibrations caused by falling leaves or the background noise of the wind or the flow of the water around rocks and other obstacles. As well as identifying the source of the vibrations, the spiders are also able to discern the distance to and direction of the source. To this end they have a range of vibration-detecting organs, including very sensitive hairs (trichobothria) on their legs and feet. Their eyes play a secondary role; experiments on related species show that touch is the main sense these spiders use to catch their prey. Their eyes are of little use for nocturnal hunting. These vibration detectors also serve to warn the spider of predators such as trout.
As soon as the vibrations reveal that there is a floundering insect within range, some fishing spiders may take direct action – they run at pace across the surface of the water and grab the insect before it extracts itself from the water and flies to safety. Some fishing spiders use silk draglines to prevent themselves from speeding past the prey.
Fishing spiders' main prey is aquatic insects, but they are opportunistic feeders and will eat anything suitable that happens within range. Dolomedes in North America have been observed catching and eating small goldfish.
Predators
The main predators of fishing spiders are birds and snakes. Dragonflies have also been observed catching young spiders. Species parasitic on the spiders include a wasp of the Pompilidae family, commonly called the Spider Wasp, that stings the spider to paralyze it before carrying it off and laying an egg in its abdomen. The larvae of the wasp hatch and proceed to eat the spider from the inside out. One escape technique the spiders use is to disappear beneath the surface tension of the water. However, some wasps, such as Anoplius depressipes, are able to be underwater for a few minutes to sting the spider and drag it out of the water.
Breeding
The males outnumber the females 3:1 suggesting a male-biased sex ratio. Mating in one North American species (D. tenebrosus) always results in the obligate death of the male, with no obvious involvement from the female.
Species
, the World Spider Catalog accepted the following 101 species:
Dolomedes actaeon Pocock, 1903 – Cameroon
Dolomedes albicomus L. Koch, 1867 – Australia (Queensland, New South Wales)
Dolomedes albicoxus Bertkau, 1880 – Brazil
Dolomedes albineus Hentz, 1845 – USA
Dolomedes alexandri Raven & Hebron, 2018 – Australia (Capital Territory, Victoria)
Dolomedes angolensis (Roewer, 1955) – Angola
Dolomedes angustivirgatus Kishida, 1936 – China, Korea, Japan
Dolomedes angustus (Thorell, 1899) – Cameroon
Dolomedes annulatus Simon, 1877 – Philippines
Dolomedes aquaticus Goyen, 1888 – New Zealand
Dolomedes batesi Pocock, 1903 – Cameroon
Dolomedes bistylus Roewer, 1955 – Congo
Dolomedes boiei (Doleschall, 1859) – Sri Lanka, Indonesia (Java)
Dolomedes briangreenei Raven & Hebron, 2018 – Australia (New South Wales, Queensland)
Dolomedes bukhkaloi Marusik, 1988 – Russia
Dolomedes chevronus Yin, 2012 – China
Dolomedes chinesus Chamberlin, 1924 – China
Dolomedes chroesus Strand, 1911 – Indonesia (Aru Is., New Guinea)
Dolomedes costatus Zhang, Zhu & Song, 2004 – China
Dolomedes crosbyi Lessert, 1928 – Congo
Dolomedes dondalei Vink & Dupérré, 2010 – New Zealand
Dolomedes elegans Taczanowski, 1874 – French Guiana
Dolomedes facetus L. Koch, 1876 – Australia, New Guinea, Samoa
Dolomedes fageli Roewer, 1955 – Congo
Dolomedes femoralis Hasselt, 1882 – Indonesia (Sumatra)
Dolomedes fernandensis Simon, 1910 – Equatorial Guinea (Bioko)
Dolomedes fimbriatus (Clerck, 1757) (type species) – Palearctic
Dolomedes flaminius L. Koch, 1867 – Australia (Queensland), New Caledonia
Dolomedes fontus Tanikawa & Miyashita, 2008 – Japan
Dolomedes furcatus Roewer, 1955 – Mozambique
Dolomedes fuscipes Roewer, 1955 – Cameroon
Dolomedes fuscus Franganillo, 1931 – Cuba
Dolomedes gertschi Carico, 1973 – USA
Dolomedes gracilipes Lessert, 1928 – Congo
Dolomedes guamuhaya Alayón, 2003 – Cuba
Dolomedes holti Carico, 1973 – Mexico
Dolomedes horishanus Kishida, 1936 – Taiwan, Japan
Dolomedes instabilis L. Koch, 1876 – Australia, Papua New Guinea
Dolomedes intermedius Giebel, 1863 – Colombia
Dolomedes japonicus Bösenberg & Strand, 1906 – China, Korea, Japan
Dolomedes kalanoro Silva & Griswold, 2013 – Madagascar
Dolomedes karijini Raven & Hebron, 2018 – Australia (Western Australia)
Dolomedes karschi Strand, 1913 – Sri Lanka
Dolomedes lafoensis Berland, 1924 – New Caledonia
Dolomedes laticeps Pocock, 1898 – Solomon Is.
Dolomedes lesserti Roewer, 1955 – Mozambique
Dolomedes lizturnerae Raven & Hebron, 2018 – Australia (Tasmania)
Dolomedes machadoi Roewer, 1955 – West Africa
Dolomedes macrops Simon, 1906 – Sudan
Dolomedes mankorlod Raven & Hebron, 2018 – Australia (Northern Territory)
Dolomedes mendigoetmopasi Barrion, 1995 – Philippines
Dolomedes minahassae Merian, 1911 – Indonesia (Sulawesi)
Dolomedes minor L. Koch, 1876 – New Zealand
Dolomedes mizhoanus Kishida, 1936 – China, Laos, Malaysia, Taiwan
Dolomedes naja Berland, 1938 – Vanuatu
Dolomedes neocaledonicus Berland, 1924 – New Caledonia
Dolomedes nigrimaculatus Song & Chen, 1991 – China, Korea
Dolomedes noukhaiva Walckenaer, 1847 – Marquesas Is.
Dolomedes ohsuditia Kishida, 1936 – Japan
Dolomedes okefinokensis Bishop, 1924 – USA
Dolomedes orion Tanikawa, 2003 – Japan
Dolomedes palmatus Zhang, Zhu & Song, 2005 – China
Dolomedes palpiger Pocock, 1903 – Cameroon
Dolomedes paroculus Simon, 1901 – Malaysia
Dolomedes pedder Raven & Hebron, 2018 – Australia (Tasmania)
Dolomedes pegasus Tanikawa, 2012 – Japan
Dolomedes petalinus Yin, 2012 – China
Dolomedes plantarius (Clerck, 1757) – Europe, Russia (Europe to South Siberia), Kazakhstan
Dolomedes pullatus Nicolet, 1849 – Chile
Dolomedes raptor Bösenberg & Strand, 1906 – Russia, China, Korea, Japan
Dolomedes raptoroides Zhang, Zhu & Song, 2004 – China
Dolomedes saganus Bösenberg & Strand, 1906 – China, Taiwan, Japan
Dolomedes schauinslandi Simon, 1899 – New Zealand
Dolomedes scriptus Hentz, 1845 – USA, Canada
Dolomedes senilis Simon, 1880 – Russia, China, Japan
Dolomedes signatus Walckenaer, 1837 – Mariana Is.
Dolomedes silvicola Tanikawa & Miyashita, 2008 – China, Japan
Dolomedes smithi Lessert, 1916 – East Africa
Dolomedes spathularis Hasselt, 1882 – Indonesia (Sumatra)
Dolomedes stilatus Karsch, 1878 – Australia
Dolomedes straeleni Roewer, 1955 – Congo
Dolomedes striatus Giebel, 1869 – USA, Canada
Dolomedes sulfureus L. Koch, 1878 – Russia, China, Korea, Japan
Dolomedes sumatranus Strand, 1906 – Indonesia (Sumatra)
Dolomedes tadzhikistanicus Andreeva, 1976 – Tajikistan
Dolomedes tenebrosus Hentz, 1844 – USA, Canada
Dolomedes titan Berland, 1924 – New Caledonia, Vanuatu
Dolomedes toldo Alayón, 2003 – Cuba
Dolomedes transfuga Pocock, 1900 – Congo
Dolomedes triton (Walckenaer, 1837) – North America, Cuba
Dolomedes upembensis (Roewer, 1955) – Congo
Dolomedes vatovae Caporiacco, 1940 – Ethiopia
Dolomedes venmani Raven & Hebron, 2018 – Australia (New South Wales, Queensland)
Dolomedes vicque Raven & Hebron, 2018 – Australia (Victoria, New South Wales, Queensland)
Dolomedes vittatus Walckenaer, 1837 – USA
Dolomedes wetarius Strand, 1911 – Indonesia
Dolomedes wollastoni Hogg, 1915 – New Guinea
Dolomedes wollemi Raven & Hebron, 2018 – Australia (New South Wales)
Dolomedes yawatai Ono, 2002 – Japan (Ryukyu Is.)
Dolomedes zatsun Tanikawa, 2003 – Japan
Dolomedes zhangjiajiensis Yin, 2012 – China
Distribution
The approximately 100 species of Dolomedes have a worldwide distribution. The largest number of species are found in Asia, with particularly high species diversity in South-east Asia, from China and Japan to New Guinea. The second largest number of species occur in tropical Africa. South America has only four species.
North America
Nine species of Dolomedes exist in North America. The six-spotted fishing spider (D. triton) lives primarily in small lakes and ponds. This spider consumes mostly water striders (pond skaters), but like all Dolomedes, it is an opportunistic ambush hunter that will eat anything that it can capture. Other species include the bog-dwelling D. striatus, and four species living by streams: D. scriptus, D. vittatus, D. gertschi and D. holti. Two North American species, D. tenebrosus and D. okefinokensis, exhibit female giganticism and/or male dwarfism, with their males being less than half the size of the females. The ninth species is the arboreal D. albineus.
Europe
Two Dolomedes species occur in Europe (excluding Russia). The Palearctic raft spider (D. fimbriatus) is widespread on the surface of bog pools and in boggy grassland. The great raft spider (D. plantarius) lives in fens, and is listed as endangered in Great Britain and is globally vulnerable.
New Zealand
Four endemic species of Dolomedes occur in New Zealand, three on the mainland and one on the Chatham Islands. Two are widespread: D. aquaticus of open riverbanks and lakeshores, and D. dondalei or New Zealand forest fishing spider (once referred to as Dolomedes III), which specialises in forested riverbanks. The largest New Zealand species, D. schauinslandi or the Rangatira spider, occurs on rodent-free islands in the Chathams where running water is rare. The fourth and most common species, D. minor, is found in scrubland, grassland, and wetlands. It mostly hunts on the ground, but is still capable of catching aquatic prey. Known as the nursery web spider, it makes white nursery webs on shrubs.
| Biology and health sciences | Spiders | Animals |
2110211 | https://en.wikipedia.org/wiki/Stari%20Most | Stari Most | Stari Most (), also known as Mostar Bridge, is a rebuilt 16th-century Ottoman bridge in the city of Mostar in Bosnia and Herzegovina. It crosses the river Neretva and connects the two parts of the city, which is named after the bridge keepers (mostari) who guarded the Stari Most during the Ottoman era. During the Croat–Bosniak War, the Army of the Republic of Bosnia and Herzegovina used the bridge as a military supply line, and the bridge was shelled by the Croatian Defence Council (HVO) and collapsed on 9 November 1993. Subsequently, the bridge was reconstructed, and it re-opened on 23 July 2004. In 2017, the ICTY deemed that the shelling was legal and that the bridge was a legitimate military target.
The Old Bridge is an exemplary piece of Balkan Islamic architecture. It was commissioned by Suleiman the Magnificent in 1557 and designed by Mimar Hayruddin, a student and apprentice of the architect Mimar Sinan.
Characteristics
The bridge spans the Neretva River in the old town of Mostar, the unofficial capital of Herzegovina. The Stari Most is hump-backed, wide and long, and dominates the river from a height of . Two fortified towers protect it: the Halebija tower on the northeast and the Tara tower on the southwest, called "the bridge keepers" (natively mostari).
Instead of foundations, the bridge has abutments of limestone linked to wing walls along the waterside cliffs. Measuring from the summer water level of , abutments are erected to a height of , from which the arch springs to its high point. The start of the arch is emphasized by a molding in height. The rise of the arch is .
History
The stone single-arch bridge is considered an exemplary piece of Balkan Islamic architecture and was commissioned by Suleiman the Magnificent in 1557. It was designed by Mimar Hayruddin, a student and apprentice of architect Mimar Sinan who built many of the Sultan's key buildings in Istanbul and around the empire.
As Mostar's economic and administrative importance grew with the growing presence of Ottoman rule, the precarious wooden suspension bridge over the Neretva gorge required replacement. The old bridge on the river "...was made of wood and hung on chains," wrote the Ottoman geographer Katip Çelebi, and it "...swayed so much that people crossing it did so in mortal fear". In 1566, Mimar Hayruddin designed the replacement bridge, which was said to have cost 300,000 Drams (silver coins) to build. The two-year construction project was supervised by Karagoz Mehmet Bey, Sultan Suleiman's son-in-law and the patron of Mostar's most important mosque complex, the Hadzi Mehmed Karadzozbeg Mosque.
Construction began in 1557 and took nine years: according to the inscription the bridge was completed in 974 AH, corresponding to the period between 19 July 1566 and 7 July 1567. Little is known of the construction of the bridge, thought to have been made from mortar made with egg whites, and all that has been preserved in writing are memories and legends and the name of the builder, Mimar Hayruddin. Charged under pain of death to construct a bridge of such unprecedented dimensions, Hayruddin reportedly prepared for his own funeral on the day the scaffolding was finally removed from the completed structure. Upon its completion it was the widest human-made arch in the world.
The 17th Century Ottoman explorer Evliya Çelebi wrote that the bridge "is like a rainbow arch soaring up to the skies, extending from one cliff to the other... I, a poor and miserable slave of Allah, have passed through 16 countries, but I have never seen such a high bridge. It is thrown from rock to rock as high as the sky."
Destruction
During the Croat–Bosniak War, the Bosniak Army of the Republic of Bosnia and Herzegovina used the Old Bridge as a military supply line. Slobodan Praljak, the commander of the Croat Defence Council ordered the destruction of the bridge which collapsed on 9 November 1993 as a result of shelling by the Bosnian Croat forces. The International Criminal Tribunal for the former Yugoslavia found it to be a legitimate military target as the opposing Army of the Republic of Bosnia and Herzegovina used it for military purposes.
The first temporary bridge on the traces of the Old Bridge was open on 30 December 1993; built in only three days by Spanish military engineers assigned to the United Nations Protection Force (UNPROFOR) mission. The temporary structure was subsequently upgraded three times, to eventually link the shores with a more secure cable-stayed bridge until the proper reconstruction of the Old Bridge.
Newspapers based in Sarajevo reported that more than 60 shells hit the bridge before it collapsed. Praljak published a document, "How the Old Bridge Was Destroyed", where he argues that there was an explosive charge or mine placed at the centre of the bridge underneath and detonated remotely, in addition to the shelling, which caused the collapse. Most historians dismiss these claims and disagree with their conclusions.
Reconstruction
After the end of the war, plans were raised to reconstruct the bridge. The World Bank, the United Nations Educational, Scientific and Cultural Organization (UNESCO), the Aga Khan Trust for Culture and the World Monuments Fund formed a coalition to oversee the reconstruction of the Stari Most and the historic city centre of Mostar. Additional funding was provided by Italy, the Netherlands, Turkey, Croatia and the Council of Europe Development Bank, as well as the Government of BiH. In October 1998, UNESCO established an international committee of experts to oversee the design and reconstruction work. It was decided to build a bridge as similar as possible to the original, using the same technology and materials.
The bridge was re-built in two phases: the first one being led by Hungarian army engineers, consisting of the lifting of submerged material for its repurpose; and the second one being the removal of the temporary bridge, a task assigned to Spanish army engineers, and the reconstruction of the Old Bridge with Ottoman construction techniques by a partnership of civil engineering companies led by the Turkish Er-Bu. Tenelia, a fine-grained limestone, sourced from local quarries was used and Hungarian army divers recovered stones from the original bridge from the river below, although most were too damaged to reuse.
Reconstruction commenced on 7 June 2001. The reconstructed bridge was inaugurated on 23 July 2004, with the cost estimated to be 15.5 million US dollars.
Diving
Stari Most diving is a traditional annual competition in diving organized every year in mid summer (end of July). It is traditional for the young men of the town to leap from the bridge into the Neretva. As the Neretva is very cold, this is a risky feat and requires skill and training, though according to TripAdvisor, tourists do dive as well. In 1968 a formal diving competition was inaugurated and held every summer. The first person to jump from the bridge since it was re-opened was Enej Kelecija.
Since 2015, Stari Most has been a tour stop in the Red Bull Cliff Diving World Series. In 2019 the diving was featured on Series 2, episode 3 of The Misadventures of Romesh Ranganathan.
In popular culture
Turkish rock band Bulutsuzluk Özlemi's 1996 song "Yaşamaya Mecbursun" (lit. 'You have to live') is about the destruction of Stari Most.Old Bridge'', a play by Papatango New Writing Prize winner Igor Memic, explores personal and historical narratives tied to the significance of the Old Bridge in Mostar. It premiered in 2021 at the Bush Theatre in London and received the Outstanding Achievement in an Affiliate Theatre award at the 2022 Olivier Awards.
| Technology | Bridges | null |
2110720 | https://en.wikipedia.org/wiki/Pencil%20fish | Pencil fish | Nannostomus (from the Greek nanos = small, and the Latin stomus = relating to the mouth) is a genus of fish belonging to the characin family Lebiasinidae. All of the species in this genus are known as pencil fish, a popular name that was initially only applied to two species in the 1920s, Nannostomus unifasciatus and Nannostomus eques, by the late 1950s however, the term would be applied to all members of the genus. Several species have become popular aquarium fish due to their attractive coloration, unique shape, and interesting demeanor.
Taxonomy
The genus Nannostomus was first described by Günther in 1872 with the type species, N. beckfordi. In 1876, Steindachner described three more species, N. unifasciatus, N. eques (pictured below), and N. trifasciatus (pictured above). In 1909, Carl H. Eigenmann described N. marginatus, N. minimus, N. erythrurus and N. harrisoni. Several of these have been popular with aquarists since the early 20th century, partly due to enthusiastic articles written about them and photographs taken by William T. Innes that were published as early as 1933. Over the years, the genus was split by subsequent authors into other genera, including Poecilobrycon and Nannobrycon. After nearly a century of debate on the subject, Dr. Stanley Howard Weitzman and Dr. J. Stanley Cobb restored earlier taxonomy and expanded upon it, unifying all species under Nannostomus in 1975. This comprehensive revision of the genus has now been widely accepted. Dr. Weitzman is also responsible for describing five of the more recently introduced species, N. marilynae, N. limatus, N. nitidus, N. britskii and N. anduzei. Twenty species are now known, most of which are also familiar to aquarists. Several other unidentified Nannostomus species have been imported over the years; many were found as bycatch with other small characins, but their taxonomic status is yet to be determined.
Species
The 20 currently recognized species in this genus are:
Nannostomus anduzei Fernández & S. H. Weitzman, 1987
Nannostomus beckfordi Günther, 1872 (golden pencil fish)
Nannostomus bifasciatus Hoedeman, 1954 (two-lined pencil fish, whiteside pencil fish)
Nannostomus britskii S. H. Weitzman, 1978 (spot stripe pencil fish)
Nannostomus digrammus (Fowler, 1913) (two-stripe pencil fish)
Nannostomus eques Steindachner, 1876 (brown pencil fish, Diptail)
Nannostomus erythrurus (C. H. Eigenmann, 1909)
Nannostomus espei (Meinken, 1956) (Espe's pencil fish, barred pencil fish, brown pencil fish)
Nannostomus grandis Zarske, 2011
Nannostomus harrisoni (C. H. Eigenmann, 1909) (Harrison's pencil fish, black stripe pencil fish)
Nannostomus limatus S. H. Weitzman, 1978 (elegant pencil fish)
Nannostomus marginatus C. H. Eigenmann, 1909 (dwarf pencil fish)
Nannostomus marilynae S. H. Weitzman & Cobb, 1975 (Marilyn's pencil fish, greenstripe pencil fish)
Nannostomus minimus C. H. Eigenmann, 1909 (least pencil fish)
Nannostomus mortenthaleri Paepke & Arendt, 2001 (coral-red pencil fish)
Nannostomus nigrotaeniatus Zarske, 2013
Nannostomus nitidus S. H. Weitzman, 1978 (shining pencil fish)
Nannostomus rubrocaudatus Zarske, 2009 (purple pencil fish)
Nannostomus trifasciatus Steindachner, 1876 (three-stripe pencil fish)
Nannostomus unifasciatus Steindachner, 1876 (one-line pencil fish)
Description
Most species are slender, pencil-shaped fish ranging in size from under ). N. marginatus, N. rubrocaudatus, and N. mortenthaleri possess shortened, blockier outlines reminiscent of pencil stubs. All but one species, Nannostomus espei, possess one to five horizontal black or brown stripes with gold or silver iridescence appearing above the primary stripe. Most also display red, orange, or maroon highlights on their fins, and many have flashes of these colors on their bodies as well. The recently described N. mortenthaleri and N. rubrocaudatus are especially vividly colored. N. espei is unique in that the horizontal stripes are only weakly present and are supplanted by five dark comma-shaped blotches. Other species take this appearance at night, but only N. espei displays the pattern permanently and in daylight. The adipose fin is present in some species and absent in others, while in certain species, such as N. eques, its presence or absence varies between individuals. All species swim horizontally, except N. unifasciatus and N. eques, which assume an oblique, 'snout-up' posture.
Sexual dimorphism occurs in the genus to varying degrees. In some species the males are more brilliantly colored, especially the fins. In other species dimorphism is far less evident. However, the anal fin is generally a reliable indicator of gender. For most species the anal fin of adult males is enlarged and elongated (as in N. espei, N. eques, etc.) and/or more colorful (as in N. harrisoni, N. marginatus, etc.). The popular aquarium species, N. trifasciatus, is an exception in this regard.
Distribution
The genus as a whole has a vast distribution in South America, from Colombia, Venezuela, and the Guyanas in the north, to the southern Amazon basin and Bolivia in the south, to Peru in the west, and Belém, Brazil, in the east. Several of the individual species have a distribution nearly as vast. As a result, many of the species are polymorphic and manifest marked color variations depending on the population. Over the years, some of these color variants have been erroneously described as separate species. Such names as 'Nannostomus ocellatus', 'N. anomalus' and 'N. auratus', among many others, are now known to be junior synonyms for the various species.
In the aquarium
To date, only two species, N. beckfordi and N. harrisoni, have been commercially raised for the aquarium trade in fisheries, mostly in Asia. All of the remaining species that find their way to home aquaria are wild-caught from South American waters. Nannostomus species thrive in home aquaria when provided with soft, moderately acidic water, low nitrate levels, and temperatures in the range of . The addition of aquatic plants is recommended, including floating varieties as they reduce the likelihood of fish jumping, a common occurrence for some of the species, especially N. espei and N. unifasciatus. They should be kept in schools of at least six, or in a community aquarium with other species of Nannostomus, other small peaceful characins, or corydoras. Aquariums with a strong water current, large tankmates, or swift-moving species should be avoided. If kept in a thickly planted single-species aquarium with the above water parameters, most species will spawn, eggs will not be eaten, and fry will be found growing in the floating plants. Baby brine shrimp, live or frozen, and other small-sized foods are required for both fry and adults. They are also avid biofilm grazers and, for most of the species, algae are under-reported staples of their diet. In most species, the males establish small territories and defend them. Their defensive actions are usually harmless, but in two species, N. mortenthaleri and N. trifasciatus, antagonistic behavior directed at conspecifics can be harmful if sufficient space and plant cover is not provided. Once acclimated to the aquarium and provided with suitable conditions, they are hardy, often living for five or more years.
| Biology and health sciences | Characiformes | Animals |
2111048 | https://en.wikipedia.org/wiki/Earthquake%20engineering | Earthquake engineering | Earthquake engineering is an interdisciplinary branch of engineering that designs and analyzes structures, such as buildings and bridges, with earthquakes in mind. Its overall goal is to make such structures more resistant to earthquakes. An earthquake (or seismic) engineer aims to construct structures that will not be damaged in minor shaking and will avoid serious damage or collapse in a major earthquake.
A properly engineered structure does not necessarily have to be extremely strong or expensive. It has to be properly designed to withstand the seismic effects while sustaining an acceptable level of damage.
Definition
Earthquake engineering is a scientific field concerned with protecting society, the natural environment, and the man-made environment from earthquakes by limiting the seismic risk to socio-economically acceptable levels. Traditionally, it has been narrowly defined as the study of the behavior of structures and geo-structures subject to seismic loading; it is considered as a subset of structural engineering, geotechnical engineering, mechanical engineering, chemical engineering, applied physics, etc. However, the tremendous costs experienced in recent earthquakes have led to an expansion of its scope to encompass disciplines from the wider field of civil engineering, mechanical engineering, nuclear engineering, and from the social sciences, especially sociology, political science, economics, and finance.
The main objectives of earthquake engineering are:
Foresee the potential consequences of strong earthquakes on urban areas and civil infrastructure.
Design, construct and maintain structures to perform at earthquake exposure up to the expectations and in compliance with building codes.
Seismic loading
Seismic loading means application of an earthquake-generated excitation on a structure (or geo-structure). It happens at contact surfaces of a structure either with the ground, with adjacent structures, or with gravity waves from tsunami. The loading that is expected at a given location on the Earth's surface is estimated by engineering seismology. It is related to the seismic hazard of the location.
Seismic performance
Earthquake or seismic performance defines a structure's ability to sustain its main functions, such as its safety and serviceability, at and after a particular earthquake exposure. A structure is normally considered safe if it does not endanger the lives and well-being of those in or around it by partially or completely collapsing. A structure may be considered serviceable if it is able to fulfill its operational functions for which it was designed.
Basic concepts of the earthquake engineering, implemented in the major building codes, assume that a building should survive a rare, very severe earthquake by sustaining significant damage but without globally collapsing. On the other hand, it should remain operational for more frequent, but less severe seismic events.
Seismic performance assessment
Engineers need to know the quantified level of the actual or anticipated seismic performance associated with the direct damage to an individual building subject to a specified ground shaking.
Such an assessment may be performed either experimentally or analytically.
Experimental assessment
Experimental evaluations are expensive tests that are typically done by placing a (scaled) model of the structure on a shake-table that simulates the earth shaking and observing its behavior. Such kinds of experiments were first performed more than a century ago. Only recently has it become possible to perform 1:1 scale testing on full structures.
Due to the costly nature of such tests, they tend to be used mainly for understanding the seismic behavior of structures, validating models and verifying analysis methods. Thus, once properly validated, computational models and numerical procedures tend to carry the major burden for the seismic performance assessment of structures.
Analytical/Numerical assessment
Seismic performance assessment or seismic structural analysis is a powerful tool of earthquake engineering which utilizes detailed modelling of the structure together with methods of structural analysis to gain a better understanding of seismic performance of building and non-building structures. The technique as a formal concept is a relatively recent development.
In general, seismic structural analysis is based on the methods of structural dynamics. For decades, the most prominent instrument of seismic analysis has been the earthquake response spectrum method which also contributed to the proposed building code's concept of today.
However, such methods are good only for linear elastic systems, being largely unable to model the structural behavior when damage (i.e., non-linearity) appears. Numerical step-by-step integration proved to be a more effective method of analysis for multi-degree-of-freedom structural systems with significant non-linearity under a transient process of ground motion excitation. Use of the finite element method is one of the most common approaches for analyzing non-linear soil structure interaction computer models.
Basically, numerical analysis is conducted in order to evaluate the seismic performance of buildings. Performance evaluations are generally carried out by using nonlinear static pushover analysis or nonlinear time-history analysis. In such analyses, it is essential to achieve accurate non-linear modeling of structural components such as beams, columns, beam-column joints, shear walls etc. Thus, experimental results play an important role in determining the modeling parameters of individual components, especially those that are subject to significant non-linear deformations. The individual components are then assembled to create a full non-linear model of the structure. Thus created models are analyzed to evaluate the performance of buildings.
The capabilities of the structural analysis software are a major consideration in the above process as they restrict the possible component models, the analysis methods available and, most importantly, the numerical robustness. The latter becomes a major consideration for structures that venture into the non-linear range and approach global or local collapse as the numerical solution becomes increasingly unstable and thus difficult to reach. There are several commercially available Finite Element Analysis software's such as CSI-SAP2000 and CSI-PERFORM-3D, MTR/SASSI, Scia Engineer-ECtools, ABAQUS, and Ansys, all of which can be used for the seismic performance evaluation of buildings. Moreover, there is research-based finite element analysis platforms such as OpenSees, MASTODON, which is based on the MOOSE Framework, RUAUMOKO and the older DRAIN-2D/3D, several of which are now open source.
Research for earthquake engineering
Research for earthquake engineering means both field and analytical investigation or experimentation intended for discovery and scientific explanation of earthquake engineering related facts, revision of conventional concepts in the light of new findings, and practical application of the developed theories.
The National Science Foundation (NSF) is the main United States government agency that supports fundamental research and education in all fields of earthquake engineering. In particular, it focuses on experimental, analytical and computational research on design and performance enhancement of structural systems.
The Earthquake Engineering Research Institute (EERI) is a leader in dissemination of earthquake engineering research related information both in the U.S. and globally.
A definitive list of earthquake engineering research related shaking tables around the world may be found in Experimental Facilities for Earthquake Engineering Simulation Worldwide. The most prominent of them is now E-Defense Shake Table in Japan.
Major U.S. research programs
NSF also supports the George E. Brown Jr. Network for Earthquake Engineering Simulation
The NSF Hazard Mitigation and Structural Engineering program (HMSE) supports research on new technologies for improving the behaviour and response of structural systems subject to earthquake hazards; fundamental research on safety and reliability of constructed systems; innovative developments in analysis and model based simulation of structural behaviour and response including soil-structure interaction; design concepts that improve structure performance and flexibility; and application of new control techniques for structural systems.
(NEES) that advances knowledge discovery and innovation for earthquakes and tsunami loss reduction of the nation's civil infrastructure and new experimental simulation techniques and instrumentation.
The NEES network features 14 geographically distributed, shared-use laboratories that support several types of experimental work: geotechnical centrifuge research, shake-table tests, large-scale structural testing, tsunami wave basin experiments, and field site research. Participating universities include: Cornell University; Lehigh University; Oregon State University; Rensselaer Polytechnic Institute; University at Buffalo, State University of New York; University of California, Berkeley; University of California, Davis; University of California, Los Angeles; University of California, San Diego; University of California, Santa Barbara; University of Illinois, Urbana-Champaign; University of Minnesota; University of Nevada, Reno; and the University of Texas, Austin.
The equipment sites (labs) and a central data repository are connected to the global earthquake engineering community via the NEEShub website. The NEES website is powered by HUBzero software developed at Purdue University for nanoHUB specifically to help the scientific community share resources and collaborate. The cyberinfrastructure, connected via Internet2, provides interactive simulation tools, a simulation tool development area, a curated central data repository, animated presentations, user support, telepresence, mechanism for uploading and sharing resources, and statistics about users and usage patterns.
This cyberinfrastructure allows researchers to: securely store, organize and share data within a standardized framework in a central location; remotely observe and participate in experiments through the use of synchronized real-time data and video; collaborate with colleagues to facilitate the planning, performance, analysis, and publication of research experiments; and conduct computational and hybrid simulations that may combine the results of multiple distributed experiments and link physical experiments with computer simulations to enable the investigation of overall system performance.
These resources jointly provide the means for collaboration and discovery to improve the seismic design and performance of civil and mechanical infrastructure systems.
Earthquake simulation
The very first earthquake simulations were performed by statically applying some horizontal inertia forces based on scaled peak ground accelerations to a mathematical model of a building. With the further development of computational technologies, static approaches began to give way to dynamic ones.
Dynamic experiments on building and non-building structures may be physical, like shake-table testing, or virtual ones. In both cases, to verify a structure's expected seismic performance, some researchers prefer to deal with so called "real time-histories" though the last cannot be "real" for a hypothetical earthquake specified by either a building code or by some particular research requirements. Therefore, there is a strong incentive to engage an earthquake simulation which is the seismic input that possesses only essential features of a real event.
Sometimes earthquake simulation is understood as a re-creation of local effects of a strong earth shaking.
Structure simulation
Theoretical or experimental evaluation of anticipated seismic performance mostly requires a structure simulation which is based on the concept of structural likeness or similarity. Similarity is some degree of analogy or resemblance between two or more objects. The notion of similarity rests either on exact or approximate repetitions of patterns in the compared items.
In general, a building model is said to have similarity with the real object if the two share geometric similarity, kinematic similarity and dynamic similarity. The most vivid and effective type of similarity is the kinematic one. Kinematic similarity exists when the paths and velocities of moving particles of a model and its prototype are similar.
The ultimate level of kinematic similarity is kinematic equivalence when, in the case of earthquake engineering, time-histories of each story lateral displacements of the model and its prototype would be the same.
Seismic vibration control
Seismic vibration control is a set of technical means aimed to mitigate seismic impacts in building and non-building structures. All seismic vibration control devices may be classified as passive, active or hybrid where:
passive control devices have no feedback capability between them, structural elements and the ground;
active control devices incorporate real-time recording instrumentation on the ground integrated with earthquake input processing equipment and actuators within the structure;
hybrid control devices have combined features of active and passive control systems.
When ground seismic waves reach up and start to penetrate a base of a building, their energy flow density, due to reflections, reduces dramatically: usually, up to 90%. However, the remaining portions of the incident waves during a major earthquake still bear a huge devastating potential.
After the seismic waves enter a superstructure, there are a number of ways to control them in order to soothe their damaging effect and improve the building's seismic performance, for instance:
to dissipate the wave energy inside a superstructure with properly engineered dampers;
to disperse the wave energy between a wider range of frequencies;
to absorb the resonant portions of the whole wave frequencies band with the help of so-called mass dampers.
Devices of the last kind, abbreviated correspondingly as TMD for the tuned (passive), as AMD for the active, and as HMD for the hybrid mass dampers, have been studied and installed in high-rise buildings, predominantly in Japan, for a quarter of a century.
However, there is quite another approach: partial suppression of the seismic energy flow into the superstructure known as seismic or base isolation.
For this, some pads are inserted into or under all major load-carrying elements in the base of the building which should substantially decouple a superstructure from its substructure resting on a shaking ground.
The first evidence of earthquake protection by using the principle of base isolation was discovered in Pasargadae, a city in ancient Persia, now Iran, and dates back to the 6th century BCE. Below, there are some samples of seismic vibration control technologies of today.
Dry-stone walls in Peru
Peru is a highly seismic land; for centuries the dry-stone construction proved to be more earthquake-resistant than using mortar.
People of Inca civilization were masters of the polished 'dry-stone walls', called ashlar, where blocks of stone were cut to fit together tightly without any mortar. The Incas were among the best stonemasons the world has ever seen and many junctions in their masonry were so perfect that even blades of grass could not fit between the stones.
The stones of the dry-stone walls built by the Incas could move slightly and resettle without the walls collapsing, a passive structural control technique employing both the principle of energy dissipation (coulomb damping) and that of suppressing resonant amplifications.
Tuned mass damper
Typically the tuned mass dampers are huge concrete blocks mounted in skyscrapers or other structures and move in opposition to the resonance frequency oscillations of the structures by means of some sort of spring mechanism.
The Taipei 101 skyscraper needs to withstand typhoon winds and earthquake tremors common in this area of Asia/Pacific. For this purpose, a steel pendulum weighing 660 metric tonnes that serves as a tuned mass damper was designed and installed atop the structure. Suspended from the 92nd to the 88th floor, the pendulum sways to decrease resonant amplifications of lateral displacements in the building caused by earthquakes and strong gusts.
Hysteretic dampers
A hysteretic damper is intended to provide better and more reliable seismic performance than that of a conventional structure by increasing the dissipation of seismic input energy. There are five major groups of hysteretic dampers used for the purpose, namely:
Fluid viscous dampers (FVDs)
Viscous Dampers have the benefit of being a supplemental damping system. They have an oval hysteretic loop and the damping is velocity dependent. While some minor maintenance is potentially required, viscous dampers generally do not need to be replaced after an earthquake. While more expensive than other damping technologies they can be used for both seismic and wind loads and are the most commonly used hysteretic damper.
Friction dampers (FDs)
Friction dampers tend to be available in two major types, linear and rotational and dissipate energy by heat. The damper operates on the principle of a coulomb damper. Depending on the design, friction dampers can experience stick-slip phenomenon and Cold welding. The main disadvantage being that friction surfaces can wear over time and for this reason they are not recommended for dissipating wind loads. When used in seismic applications wear is not a problem and there is no required maintenance. They have a rectangular hysteretic loop and as long as the building is sufficiently elastic they tend to settle back to their original positions after an earthquake.
Metallic yielding dampers (MYDs)
Metallic yielding dampers, as the name implies, yield in order to absorb the earthquake's energy. This type of damper absorbs a large amount of energy however they must be replaced after an earthquake and may prevent the building from settling back to its original position.
Viscoelastic dampers (VEDs)
Viscoelastic dampers are useful in that they can be used for both wind and seismic applications, they are usually limited to small displacements. There is some concern as to the reliability of the technology as some brands have been banned from use in buildings in the United States.
Straddling pendulum dampers (swing)
Base isolation
Base isolation seeks to prevent the kinetic energy of the earthquake from being transferred into elastic energy in the building. These technologies do so by isolating the structure from the ground, thus enabling them to move somewhat independently. The degree to which the energy is transferred into the structure and how the energy is dissipated will vary depending on the technology used.
Lead rubber bearing
Lead rubber bearing or LRB is a type of base isolation employing a heavy damping. It was invented by Bill Robinson, a New Zealander.
Heavy damping mechanism incorporated in vibration control technologies and, particularly, in base isolation devices, is often considered a valuable source of suppressing vibrations thus enhancing a building's seismic performance. However, for the rather pliant systems such as base isolated structures, with a relatively low bearing stiffness but with a high damping, the so-called "damping force" may turn out the main pushing force at a strong earthquake. The video shows a Lead Rubber Bearing being tested at the UCSD Caltrans-SRMD facility. The bearing is made of rubber with a lead core. It was a uniaxial test in which the bearing was also under a full structure load. Many buildings and bridges, both in New Zealand and elsewhere, are protected with lead dampers and lead and rubber bearings. Te Papa Tongarewa, the national museum of New Zealand, and the New Zealand Parliament Buildings have been fitted with the bearings. Both are in Wellington which sits on an active fault.
Springs-with-damper base isolator
Springs-with-damper base isolator installed under a three-story town-house, Santa Monica, California is shown on the photo taken prior to the 1994 Northridge earthquake exposure. It is a base isolation device conceptually similar to Lead Rubber Bearing.
One of two three-story town-houses like this, which was well instrumented for recording of both vertical and horizontal accelerations on its floors and the ground, has survived a severe shaking during the Northridge earthquake and left valuable recorded information for further study.
Simple roller bearing
Simple roller bearing is a base isolation device which is intended for protection of various building and non-building structures against potentially damaging lateral impacts of strong earthquakes.
This metallic bearing support may be adapted, with certain precautions, as a seismic isolator to skyscrapers and buildings on soft ground. Recently, it has been employed under the name of metallic roller bearing for a housing complex (17 stories) in Tokyo, Japan.
Friction pendulum bearing (FPB) is another name of friction pendulum system (FPS). It is based on three pillars:
articulated friction slider;
spherical concave sliding surface;
enclosing cylinder for lateral displacement restraint.
Snapshot with the link to video clip of a shake-table testing of FPB system supporting a rigid building model is presented at the right.
Seismic design
Seismic design is based on authorized engineering procedures, principles and criteria meant to design or retrofit structures subject to earthquake exposure. Those criteria are only consistent with the contemporary state of the knowledge about earthquake engineering structures. Therefore, a building design which exactly follows seismic code regulations does not guarantee safety against collapse or serious damage.
The price of poor seismic design may be enormous. Nevertheless, seismic design has always been a trial and error process whether it was based on physical laws or on empirical knowledge of the structural performance of different shapes and materials.
To practice seismic design, seismic analysis or seismic evaluation of new and existing civil engineering projects, an engineer should, normally, pass examination on Seismic Principles which, in the State of California, include:
Seismic Data and Seismic Design Criteria
Seismic Characteristics of Engineered Systems
Seismic Forces
Seismic Analysis Procedures
Seismic Detailing and Construction Quality Control
To build up complex structural systems, seismic design largely uses the same relatively small number of basic structural elements (to say nothing of vibration control devices) as any non-seismic design project.
Normally, according to building codes, structures are designed to "withstand" the largest earthquake of a certain probability that is likely to occur at their location. This means the loss of life should be minimized by preventing collapse of the buildings.
Seismic design is carried out by understanding the possible failure modes of a structure and providing the structure with appropriate strength, stiffness, ductility, and configuration to ensure those modes cannot occur.
Seismic design requirements
Seismic design requirements depend on the type of the structure, locality of the project and its authorities which stipulate applicable seismic design codes and criteria. For instance, California Department of Transportation's requirements called The Seismic Design Criteria (SDC) and aimed at the design of new bridges in California incorporate an innovative seismic performance-based approach.
The most significant feature in the SDC design philosophy is a shift from a force-based assessment of seismic demand to a displacement-based assessment of demand and capacity. Thus, the newly adopted displacement approach is based on comparing the elastic displacement demand to the inelastic displacement capacity of the primary structural components while ensuring a minimum level of inelastic capacity at all potential plastic hinge locations.
In addition to the designed structure itself, seismic design requirements may include a ground stabilization underneath the structure: sometimes, heavily shaken ground breaks up which leads to collapse of the structure sitting upon it.
The following topics should be of primary concerns: liquefaction; dynamic lateral earth pressures on retaining walls; seismic slope stability; earthquake-induced settlement.
Nuclear facilities should not jeopardise their safety in case of earthquakes or other hostile external events. Therefore, their seismic design is based on criteria far more stringent than those applying to non-nuclear facilities. The Fukushima I nuclear accidents and damage to other nuclear facilities that followed the 2011 Tōhoku earthquake and tsunami have, however, drawn attention to ongoing concerns over Japanese nuclear seismic design standards and caused many other governments to re-evaluate their nuclear programs. Doubt has also been expressed over the seismic evaluation and design of certain other plants, including the Fessenheim Nuclear Power Plant in France.
Failure modes
Failure mode is the manner by which an earthquake induced failure is observed. It, generally, describes the way the failure occurs. Though costly and time-consuming, learning from each real earthquake failure remains a routine recipe for advancement in seismic design methods. Below, some typical modes of earthquake-generated failures are presented.
The lack of reinforcement coupled with poor mortar and inadequate roof-to-wall ties can result in substantial damage to an unreinforced masonry building. Severely cracked or leaning walls are some of the most common earthquake damage. Also hazardous is the damage that may occur between the walls and roof or floor diaphragms. Separation between the framing and the walls can jeopardize the vertical support of roof and floor systems.
Soft story effect. Absence of adequate stiffness on the ground level caused damage to this structure. A close examination of the image reveals that the rough board siding, once covered by a brick veneer, has been completely dismantled from the studwall. Only the rigidity of the floor above combined with the support on the two hidden sides by continuous walls, not penetrated with large doors as on the street sides, is preventing full collapse of the structure.
Soil liquefaction. In the cases where the soil consists of loose granular deposited materials with the tendency to develop excessive hydrostatic pore water pressure of sufficient magnitude and compact, liquefaction of those loose saturated deposits may result in non-uniform settlements and tilting of structures. This caused major damage to thousands of buildings in Niigata, Japan during the 1964 earthquake.
Landslide rock fall. A landslide is a geological phenomenon which includes a wide range of ground movement, including rock falls. Typically, the action of gravity is the primary driving force for a landslide to occur though in this case there was another contributing factor which affected the original slope stability: the landslide required an earthquake trigger before being released.
Pounding against adjacent building. This is a photograph of the collapsed five-story tower, St. Joseph's Seminary, Los Altos, California which resulted in one fatality. During Loma Prieta earthquake, the tower pounded against the independently vibrating adjacent building behind. A possibility of pounding depends on both buildings' lateral displacements which should be accurately estimated and accounted for.
At Northridge earthquake, the Kaiser Permanente concrete frame office building had joints completely shattered, revealing inadequate confinement steel, which resulted in the second story collapse. In the transverse direction, composite end shear walls, consisting of two wythes of brick and a layer of shotcrete that carried the lateral load, peeled apart because of inadequate through-ties and failed.
Improper construction site on a foothill.
Poor detailing of the reinforcement (lack of concrete confinement in the columns and at the beam-column joints, inadequate splice length).
Seismically weak soft story at the first floor.
Long cantilevers with heavy dead load.
Sliding off foundations effect of a relatively rigid residential building structure during 1987 Whittier Narrows earthquake. The magnitude 5.9 earthquake pounded the Garvey West Apartment building in Monterey Park, California and shifted its superstructure about 10 inches to the east on its foundation.
If a superstructure is not mounted on a base isolation system, its shifting on the basement should be prevented.
Reinforced concrete column burst at Northridge earthquake due to insufficient shear reinforcement mode which allows main reinforcement to buckle outwards. The deck unseated at the hinge and failed in shear. As a result, the La Cienega-Venice underpass section of the 10 Freeway collapsed.
Loma Prieta earthquake: side view of reinforced concrete support-columns failure which triggered the upper deck collapse onto the lower deck of the two-level Cypress viaduct of Interstate Highway 880, Oakland, CA.
Retaining wall failure at Loma Prieta earthquake in Santa Cruz Mountains area: prominent northwest-trending extensional cracks up to 12 cm (4.7 in) wide in the concrete spillway to Austrian Dam, the north abutment.
Ground shaking triggered soil liquefaction in a subsurface layer of sand, producing differential lateral and vertical movement in an overlying carapace of unliquefied sand and silt. This mode of ground failure, termed lateral spreading, is a principal cause of liquefaction-related earthquake damage.
Severely damaged building of Agriculture Development Bank of China after 2008 Sichuan earthquake: most of the beams and pier columns are sheared. Large diagonal cracks in masonry and veneer are due to in-plane loads while abrupt settlement of the right end of the building should be attributed to a landfill which may be hazardous even without any earthquake.
Twofold tsunami impact: sea waves hydraulic pressure and inundation. Thus, the Indian Ocean earthquake of December 26, 2004, with the epicenter off the west coast of Sumatra, Indonesia, triggered a series of devastating tsunamis, killing more than 230,000 people in eleven countries by inundating surrounding coastal communities with huge waves up to 30 meters (100 feet) high.
Earthquake-resistant construction
Earthquake construction means implementation of seismic design to enable building and non-building structures to live through the anticipated earthquake exposure up to the expectations and in compliance with the applicable building codes.
Design and construction are intimately related. To achieve a good workmanship, detailing of the members and their connections should be as simple as possible. As any construction in general, earthquake construction is a process that consists of the building, retrofitting or assembling of infrastructure given the construction materials available.
The destabilizing action of an earthquake on constructions may be direct (seismic motion of the ground) or indirect (earthquake-induced landslides, soil liquefaction and waves of tsunami).
A structure might have all the appearances of stability, yet offer nothing but danger when an earthquake occurs. The crucial fact is that, for safety, earthquake-resistant construction techniques are as important as quality control and using correct materials. Earthquake contractor should be registered in the state/province/country of the project location (depending on local regulations), bonded and insured .
To minimize possible losses, construction process should be organized with keeping in mind that earthquake may strike any time prior to the end of construction.
Each construction project requires a qualified team of professionals who understand the basic features of seismic performance of different structures as well as construction management.
Adobe structures
Around thirty percent of the world's population lives or works in earth-made construction. Adobe type of mud bricks is one of the oldest and most widely used building materials. The use of adobe is very common in some of the world's most hazard-prone regions, traditionally across Latin America, Africa, Indian subcontinent and other parts of Asia, Middle East and Southern Europe.
Adobe buildings are considered very vulnerable at strong quakes. However, multiple ways of seismic strengthening of new and existing adobe buildings are available.
Key factors for the improved seismic performance of
adobe construction are:
Quality of construction.
Compact, box-type layout.
Seismic reinforcement.
Limestone and sandstone structures
Limestone is very common in architecture, especially in North America and Europe. Many landmarks across the world are made of limestone. Many medieval churches and castles in Europe are made of limestone and sandstone masonry. They are the long-lasting materials but their rather heavy weight is not beneficial for adequate seismic performance.
Application of modern technology to seismic retrofitting can enhance the survivability of unreinforced masonry structures. As an example, from 1973 to 1989, the Salt Lake City and County Building in Utah was exhaustively renovated and repaired with an emphasis on preserving historical accuracy in appearance. This was done in concert with a seismic upgrade that placed the weak sandstone structure on base isolation foundation to better protect it from earthquake damage.
Timber frame structures
Timber framing dates back thousands of years, and has been used in many parts of the world during various periods such as ancient Japan, Europe and medieval England in localities where timber was in good supply and building stone and the skills to work it were not.
The use of timber framing in buildings provides their complete skeletal framing which offers some structural benefits as the timber frame, if properly engineered, lends itself to better seismic survivability.
Light-frame structures
Light-frame structures usually gain seismic resistance from rigid plywood shear walls and wood structural panel diaphragms. Special provisions for seismic load-resisting systems for all engineered wood structures requires consideration of diaphragm ratios, horizontal and vertical diaphragm shears, and connector/fastener values. In addition, collectors, or drag struts, to distribute shear along a diaphragm length are required.
Reinforced masonry structures
A construction system where steel reinforcement is embedded in the mortar joints of masonry or placed in holes and that are filled with concrete or grout is called reinforced masonry. There are various practices and techniques to reinforce masonry. The most common type is the reinforced hollow unit masonry.
To achieve a ductile behavior in masonry, it is necessary that the shear strength of the wall is greater than the flexural strength. The effectiveness of both vertical and horizontal reinforcements depends on the type and quality of the masonry units and mortar.
The devastating 1933 Long Beach earthquake revealed that masonry is prone to earthquake damage, which led to the California State Code making masonry reinforcement mandatory across California.
Reinforced concrete structures
Reinforced concrete is concrete in which steel reinforcement bars (rebars) or fibers have been incorporated to strengthen a material that would otherwise be brittle. It can be used to produce beams, columns, floors or bridges.
Prestressed concrete is a kind of reinforced concrete used for overcoming concrete's natural weakness in tension. It can be applied to beams, floors or bridges with a longer span than is practical with ordinary reinforced concrete. Prestressing tendons (generally of high tensile steel cable or rods) are used to provide a clamping load which produces a compressive stress that offsets the tensile stress that the concrete compression member would, otherwise, experience due to a bending load.
To prevent catastrophic collapse in response earth shaking (in the interest of life safety), a traditional reinforced concrete frame should have ductile joints. Depending upon the methods used and the imposed seismic forces, such buildings may be immediately usable, require extensive repair, or may have to be demolished.
Prestressed structures
Prestressed structure is the one whose overall integrity, stability and security depend, primarily, on a prestressing. Prestressing means the intentional creation of permanent stresses in a structure for the purpose of improving its performance under various service conditions.
There are the following basic types of prestressing:
Pre-compression (mostly, with the own weight of a structure)
Pretensioning with high-strength embedded tendons
Post-tensioning with high-strength bonded or unbonded tendons
Today, the concept of prestressed structure is widely engaged in design of buildings, underground structures, TV towers, power stations, floating storage and offshore facilities, nuclear reactor vessels, and numerous kinds of bridge systems.
A beneficial idea of prestressing was, apparently, familiar to the ancient Roman architects; look, e.g., at the tall attic wall of Colosseum working as a stabilizing device for the wall piers beneath.
Steel structures
Steel structures are considered mostly earthquake resistant but some failures have occurred. A great number of welded steel moment-resisting frame buildings, which looked earthquake-proof, surprisingly experienced brittle behavior and were hazardously damaged in the 1994 Northridge earthquake. After that, the Federal Emergency Management Agency (FEMA) initiated development of repair techniques and new design approaches to minimize damage to steel moment frame buildings in future earthquakes.
For structural steel seismic design based on Load and Resistance Factor Design (LRFD) approach, it is very important to assess ability of a structure to develop and maintain its bearing resistance in the inelastic range. A measure of this ability is ductility, which may be observed in a material itself, in a structural element, or to a whole structure.
As a consequence of Northridge earthquake experience, the American Institute of Steel Construction has introduced AISC 358 "Pre-Qualified Connections for Special and intermediate Steel Moment Frames." The AISC Seismic Design Provisions require that all Steel Moment Resisting Frames employ either connections contained in AISC 358, or the use of connections that have been subjected to pre-qualifying cyclic testing.
Prediction of earthquake losses
Earthquake loss estimation is usually defined as a Damage Ratio (DR) which is a ratio of the earthquake damage repair cost to the total value of a building. Probable Maximum Loss (PML) is a common term used for earthquake loss estimation, but it lacks a precise definition. In 1999, ASTM E2026 'Standard Guide for the Estimation of Building Damageability in Earthquakes' was produced in order to standardize the nomenclature for seismic loss estimation, as well as establish guidelines as to the review process and qualifications of the reviewer.
Earthquake loss estimations are also referred to as Seismic Risk Assessments. The risk assessment process generally involves determining the probability of various ground motions coupled with the vulnerability or damage of the building under those ground motions. The results are defined as a percent of building replacement value.
| Technology | Disciplines | null |
2111904 | https://en.wikipedia.org/wiki/Alpheidae | Alpheidae | Alpheidae (also known as the snapping shrimp, pistol shrimp or alpheid shrimp) is a family within the infraorder caridea characterized by having asymmetrical claws, the larger of which is typically capable of producing a loud snapping sound.
The family is diverse and worldwide in distribution, consisting of about 1,119 species within 38 or more genera. The two most prominent genera are Alpheus and Synalpheus, with species numbering well over 330 and 160, respectively. Most snapping shrimp dig burrows and are common inhabitants of coral reefs, submerged seagrass flats, and oyster reefs. While most genera and species are found in tropical and temperate coastal and marine waters, Betaeus inhabits cold seas and Potamalpheops has a cosmopolitan distribution including being found in freshwater caves in Mexico.
When in colonies, the snapping shrimp can interfere with sonar and underwater communication. The shrimp are considered a major source of sound in the ocean.
Characteristics
Snapping shrimp grow to in length.
Its disproportionately large claw, larger than half the shrimp's body, is a dismorphic addition to the arsenal of the shrimp. The claw can be on either arm of the body, and, unlike most shrimp claws, does not have typical pincers at the end. Rather, it has a pistol-like feature made of two parts. A joint allows the "hammer" part to move backward into a right-angled position. When released, it snaps into the other part of the claw, emitting an enormously powerful wave of bubbles capable of stunning larger fish and breaking small glass jars.
The claw snaps to create a cavitation bubble that generates acoustic pressures of up to at a distance of 4 cm from the claw. As it ejects from the claw, the bubble reaches speeds of . The pressure is high enough to kill small fish. It corresponds to a peak pressure level of 218 decibels relative to one micropascal (dB re 1 μPa), equivalent to a zero to peak source level of 190 dB re 1 μPa m. Au and Banks measured peak to peak source levels between 185 and 190 dB re 1 μPa m, depending on the size of the claw. Similar values are reported by Ferguson and Cleary. The duration of the click is less than 1 millisecond.
The snap can also produce sonoluminescence from the collapsing cavitation bubble. As it collapses, the cavitation bubble emits a short flash of light with a broad spectrum. If the light were of thermal origin it would require a temperature of the emitter of over . In comparison, the surface temperature of the Sun is estimated to be around . The light is of lower intensity than the light produced by typical sonoluminescence and is not visible to the naked eye. It is most likely a by-product of the shock wave with no biological significance. However, it was the first known instance of an animal producing light by this effect. It has subsequently been discovered that another group of crustaceans, the mantis shrimp, contains species whose club-like forelimbs can strike so quickly and with such force as to induce sonoluminescent cavitation bubbles upon impact.
The snapping is used for hunting (hence the alternative name "pistol shrimp"), as well as for communication. When hunting, the shrimp usually lies in an obscured spot, such as a burrow. The shrimp then extends its antennae outwards to determine if any fish are passing by. Once it feels movement, the shrimp inches out of its hiding place, pulls back its claw, and releases a "shot" which stuns the prey; the shrimp then pulls it to the burrow and feeds on it.
Snapping shrimp have the ability to reverse claws. When the snapping claw is lost, the missing limb will regenerate into a smaller claw and the original smaller appendage will grow into a new snapping claw. Laboratory research has shown that severing the nerve of the snapping claw induces the conversion of the smaller limb into a second snapping claw. The reversal of claw asymmetry in snapping shrimp is thought to be unique in nature.
The snapping shrimp competes with much larger animals such as the sperm whale and beluga whale for the title of loudest animal in the sea. When in colonies, the snapping shrimp can interfere with sonar and underwater communication. The shrimp are a major source of noise in the ocean and can interfere with anti-submarine warfare.
Ecology
Some snapping shrimp species share burrows with goby fish in a mutualistic symbiotic relationship. The burrow is built and tended by the pistol shrimp, and the goby provides protection by watching out for danger. When both are out of the burrow, the shrimp maintains contact with the goby using its antennae. The goby, having better vision, alerts the shrimp of danger using a characteristic tail movement, and then both retreat into the safety of the shared burrow. This association has been observed in species that inhabit coral reef habitats.
Eusocial behavior has been discovered in the genus Synalpheus. The species Synalpheus regalis lives inside sponges in colonies that can number over 300. All of them are the offspring of a single large female, the queen, and possibly a single male. The offspring are divided into workers who care for the young and predominantly male soldiers who protect the colony with their huge claws.
The snapping shrimp species will retain the same mate after copulation, making them monogamous. Most females of the Alpheidae species are susceptible to mating. Young females become receptive to males either just before (premolt stage) or after the puberty molt, making them physiologically mature and morphologically able to carry the egg mass. Male presence during the molt is beneficial for the female, as searching for a male during her soft-bodied receptive phase would put her at mortal risk. Mates have more success with partners having greater body mass. The larger shrimp are most successful. These animals practice mate guarding, leading to a decline in mate competition, as well as bonding of partners. The male and female will defend their shelter to protect both territory and young. Larva develop in three stages: The nauplius larvae, zoea, and post larval stages.
Genera
More than 620 species are currently recognised in the family Alpheidae, distributed among 52 genera. The largest of these are Alpheus, with 336 species, and Synalpheus, with 168 species. The following genera are recognised in the family Alpheidae:
Acanthanas
Alpheopsis
Alpheus
Amphibetaeus
Arete
Aretopsis
Athanas
Athanopsis
Automate
Bannereus
Batella
Bermudacaris
Betaeopsis
Betaeus
Bruceopsis
Caligoneus
Coronalpheus
Coutieralpheus
Crosnierocaris
Deioneus
Fenneralpheus
Harperalpheus
Jengalpheops
Leptalpheus
Leptathanas
Leslibetaeus
Metabetaeus
Metalpheus
Mohocaris
Nennalpheus
Notalpheus
Oligorostra
Oligosella
Orygmalpheus
Pachelpheus
Parabetaeus
Pomagnathus
Potamalpheops
Prionalpheus
Pseudalpheopsis
Pseudathanas
Pterocaris
Racilius
Richalpheus
Rugathanas
Salmoneus
Stenalpheops
Synalpheus
Thuylamea
Triacanthoneus
Vexillipar
Yagerocaris
| Biology and health sciences | Shrimps and prawns | Animals |
23096495 | https://en.wikipedia.org/wiki/Air%20separation | Air separation | An air separation plant separates atmospheric air into its primary components, typically nitrogen and oxygen, and sometimes also argon and other rare inert gases.
The most common method for air separation is fractional distillation. Cryogenic air separation units (ASUs) are built to provide nitrogen or oxygen and often co-produce argon. Other methods such as membrane, pressure swing adsorption (PSA) and vacuum pressure swing adsorption (VPSA) are commercially used to separate a single component from ordinary air. High purity oxygen, nitrogen, and argon, used for semiconductor device fabrication, require cryogenic distillation. Similarly, the only viable source of the rare gases neon, krypton, xenon is the distillation of air using at least two distillation columns. Helium is also recovered in advanced air separation processes.
Cryogenic distillation process
Pure gases can be separated from air by first cooling it until it liquefies, then selectively distilling the components at their various boiling temperatures. The process can produce high purity gases but is energy-intensive. This process was pioneered by Carl von Linde in the early 20th century and is still used today to produce high purity gases. He developed it in the year 1895; the process remained purely academic for seven years before it was used in industrial applications for the first time (1902).
The cryogenic separation process requires a very tight integration of heat exchangers and separation columns to obtain a good efficiency and all the energy for refrigeration is provided by the compression of the air at the inlet of the unit.
To achieve the low distillation temperatures, an air separation unit requires a refrigeration cycle that operates by means of the Joule–Thomson effect.
The separated products are sometimes supplied by pipeline to large industrial users near the production plant. Long distance transportation of products is by shipping liquid product for large quantities or as dewar flasks or gas cylinders for small quantities.
Non-cryogenic processes
Pressure swing adsorption provides separation of oxygen or nitrogen from air without liquefaction. The process operates around ambient temperature; a zeolite (molecular sponge) is exposed to high pressure air, then the air is released and an adsorbed film of the desired gas is released. The size of compressor is much reduced over a liquefaction plant, and portable oxygen concentrators are made in this manner to provide oxygen-enriched air for medical purposes. Vacuum swing adsorption is a similar process; the product gas is evolved from the zeolite at sub-atmospheric pressure.
Membrane technologies can provide alternate, lower-energy approaches to air separation. For example, a number of approaches are being explored for oxygen generation. Polymeric membranes operating at ambient or warm temperatures, for example, may be able to produce oxygen-enriched air (25-50% oxygen). Ceramic membranes can provide high-purity oxygen (90% or more) but require higher temperatures (800-900 deg C) to operate. These ceramic membranes include ion transport membranes (ITM) and oxygen transport membranes (OTM). Air Products and Chemicals Inc and Praxair are developing flat ITM and tubular OTM systems.
Membrane gas separation is used to provide oxygen-poor and nitrogen-rich gases instead of air to fill the fuel tanks of jet liners, thus greatly reducing the chances of accidental fires and explosions. Conversely, membrane gas separation is currently used to provide oxygen-enriched air to pilots flying at great altitudes in aircraft without pressurized cabins.
Oxygen-enriched air can be obtained exploiting the different solubility of oxygen and nitrogen. Oxygen is more soluble than nitrogen in water, so if air is degassed from water, a stream of 35% oxygen can be obtained.
Applications
Rocketry
Liquid oxygen for companies such as SpaceX.
Medical
Pure oxygen is delivered to large hospitals for use with patients.
Steel
In steelmaking, oxygen is required for the basic oxygen steelmaking process. Modern basic oxygen steelmaking uses almost two tons of oxygen per ton of steel.
Ammonia
Nitrogen used in the Haber process to make ammonia.
Coal gas
Large amounts of oxygen are required for coal gasification projects; cryogenic plants producing 3000 tons/day are found in some projects.
Inert gas
Inerting with nitrogen storage tanks of ships and tanks for petroleum products, or for protecting edible oil products from oxidation.
| Physical sciences | Phase separations | Chemistry |
23099899 | https://en.wikipedia.org/wiki/Environmental%20issues | Environmental issues | Environmental issues are disruptions in the usual function of ecosystems. Further, these issues can be caused by humans (human impact on the environment) or they can be natural. These issues are considered serious when the ecosystem cannot recover in the present situation, and catastrophic if the ecosystem is projected to certainly collapse.
Environmental protection is the practice of protecting the natural environment on the individual, organizational or governmental levels, for the benefit of both the environment and humans. Environmentalism is a social and environmental movement that addresses environmental issues through advocacy, legislation education, and activism.
Environment destruction caused by humans is a global, ongoing problem. Water pollution also cause problems to marine life. Some scholars believe that the projected peak global population of roughly 9-10 billion people could live sustainably within the earth's ecosystems if humans worked to live sustainably within planetary boundaries. The bulk of environmental impacts are caused by excessive consumption of industrial goods by the world's wealthiest populations. The UN Environmental Program, in its "Making Peace With Nature" Report in 2021, found addressing key planetary crises, like pollution, climate change and biodiversity loss, was achievable if parties work to address the Sustainable Development Goals.
Types
Major current environmental issues may include climate change, pollution, environmental degradation, and resource depletion. The conservation movement lobbies for protection of endangered species and protection of any ecologically valuable
natural areas, genetically modified foods and global warming. The UN system has adopted international frameworks for environmental issues in three key issues, which has been encoded as the "triple planetary crises": climate change, pollution, and biodiversity loss.
Human impact
Degradation
Conflict
Costs
Action
Justice
The 2023 IPCC report highlighted the disproportionate effects of climate change on vulnerable populations. The report's findings make it clear that every increment of global warming exacerbates challenges such as extreme heatwaves, heavy rainfall, and other weather extremes, which in turn amplify risks for human health and ecosystems. With nearly half of the world's population residing in regions highly susceptible to climate change, the urgency for global actions that are both rapid and sustained is underscored. The importance of integrating diverse knowledge systems, including scientific, Indigenous, and local knowledge, into climate action is highlighted as a means to foster inclusive solutions that address the complexities of climate impacts across different communities.
In addition, the report points out the critical gap in adaptation finance, noting that developing countries require significantly more resources to effectively adapt to climate challenges than what is currently available. This financial disparity raises questions about the global commitment to equitable climate action and underscores the need for a substantial increase in support and resources. The IPCC's analysis suggests that with adequate financial investment and international cooperation, it is possible to embark on a pathway towards resilience and sustainability that benefits all sections of society.
Law
Assessment
Movement
Organizations
Environmental issues are addressed at a regional, national or international level by government organizations.
The largest international agency, set up in 1972, is the United Nations Environment Programme. The International Union for Conservation of Nature brings together 83 states, 108 government agencies, 766 Non-governmental organizations and 81 international organizations and about 10,000 experts, scientists from countries around the world. International non-governmental organizations include Greenpeace, Friends of the Earth and World Wide Fund for Nature. Governments enact environmental policy and enforce environmental law and this is done to differing degrees around the world.
Film and television
There are an increasing number of films being produced on environmental issues, especially on climate change and global warming. Al Gore's 2006 film An Inconvenient Truth gained commercial success and a high media profile.
| Physical sciences | Earth science basics: General | Earth science |
23105042 | https://en.wikipedia.org/wiki/Primordial%20nuclide | Primordial nuclide | In geochemistry, geophysics and nuclear physics, primordial nuclides, also known as primordial isotopes, are nuclides found on Earth that have existed in their current form since before Earth was formed. Primordial nuclides were present in the interstellar medium from which the solar system was formed, and were formed in, or after, the Big Bang, by nucleosynthesis in stars and supernovae followed by mass ejection, by cosmic ray spallation, and potentially from other processes. They are the stable nuclides plus the long-lived fraction of radionuclides surviving in the primordial solar nebula through planet accretion until the present; 286 such nuclides are known.
Stability
All of the known 251 stable nuclides, plus another 35 nuclides that have half-lives long enough to have survived from the formation of the Earth, occur as primordial nuclides. These 35 primordial radionuclides represent isotopes of 28 separate elements.
Cadmium, tellurium, xenon, neodymium, samarium, osmium, and uranium each have two primordial radioisotopes (, ; , ; , ; , ; , ; , ; and , ).
Because the age of the Earth is (4.6 billion years), the half-life of the given nuclides must be greater than about (100 million years) for practical considerations. For example, for a nuclide with half-life (60 million years), this means 77 half-lives have elapsed, meaning that for each mole () of that nuclide being present at the formation of Earth, only 4 atoms remain today.
The seven shortest-lived primordial nuclides (i.e., the nuclides with the shortest half-lives) to have been experimentally verified are (), (), (), (), (), (), and ().
These are the seven nuclides with half-lives comparable to, or somewhat less than, the estimated age of the universe. (87Rb, 187Re, 176Lu, and 232Th have half-lives somewhat longer than the age of the universe.) For a complete list of the 35 known primordial radionuclides, including the next 28 with half-lives much longer than the age of the universe, see the complete list below. For practical purposes, nuclides with half-lives much longer than the age of the universe may be treated as if they were stable. 87Rb, 187Re, 176Lu, 232Th, and 238U have half-lives long enough that their decay is limited over geological time scales; 40K and 235U have shorter half-lives and are hence severely depleted, but are still long-lived enough to persist significantly in nature.
The longest-lived isotope not proven to be primordial is , which has a half-life of , followed by () and (). 244Pu was reported to exist in nature as a primordial nuclide in 1971, but this detection could not be confirmed by further studies in 2012 and 2022.
Taking into account that all these nuclides must exist for at least , 146Sm must survive 45 half-lives (and hence be reduced by 245 ≈ ), 244Pu must survive 57 (and be reduced by a factor of 257 ≈ ), and 92Nb must survive 130 (and be reduced by 2130 ≈ ). Mathematically, considering the likely initial abundances of these nuclides, primordial 146Sm and 244Pu should persist somewhere within the Earth today, even if they are not identifiable in the relatively minor portion of the Earth's crust available to human assays, while 92Nb and all shorter-lived nuclides should not. Nuclides such as 92Nb that were present in the primordial solar nebula but have long since decayed away completely are termed extinct radionuclides if they have no other means of being regenerated. As for 244Pu, calculations suggest that as of 2022, sensitivity limits were about one order of magnitude away from detecting it as a primordial nuclide.
Because primordial chemical elements often consist of more than one primordial isotope, there are only 83 distinct primordial chemical elements. Of these, 80 have at least one observationally stable isotope and three additional primordial elements have only radioactive isotopes (bismuth, thorium, and uranium).
Naturally occurring nuclides that are not primordial
Some unstable isotopes which occur naturally (such as , , and ) are not primordial, as they must be constantly regenerated. This occurs by cosmic radiation (in the case of cosmogenic nuclides such as and ), or (rarely) by such processes as geonuclear transmutation (neutron capture of uranium in the case of and ). Other examples of common naturally occurring but non-primordial nuclides are isotopes of radon, polonium, and radium, which are all radiogenic nuclide daughters of uranium decay and are found in uranium ores. The stable argon isotope 40Ar is actually more common as a radiogenic nuclide than as a primordial nuclide, forming almost 1% of the Earth's atmosphere, which is regenerated by the beta decay of the extremely long-lived radioactive primordial isotope 40K, whose half-life is on the order of a billion years and thus has been generating argon since early in the Earth's existence. (Primordial argon was dominated by the alpha process nuclide 36Ar, which is significantly rarer than 40Ar on Earth.)
A similar radiogenic series is derived from the long-lived radioactive primordial nuclide 232Th. These nuclides are described as geogenic, meaning that they are decay or fission products of uranium or other actinides in subsurface rocks. All such nuclides have shorter half-lives than their parent radioactive primordial nuclides. Some other geogenic nuclides do not occur in the decay chains of 232Th, 235U, or 238U but can still fleetingly occur naturally as products of the spontaneous fission of one of these three long-lived nuclides, such as 126Sn, which makes up about 10−14 of all natural tin. Another, 99Tc, has also been detected. There are five other long-lived fission products known.
Primordial elements
A primordial element is a chemical element with at least one primordial nuclide. There are 251 stable primordial nuclides and 35 radioactive primordial nuclides, but only 80 primordial stable elements—hydrogen through lead, atomic numbers 1 to 82, except for technetium (43) and promethium (61)—and three radioactive primordial elements—bismuth (83), thorium (90), and uranium (92). If plutonium (94) turns out to be primordial (specifically, the long-lived isotope Pu), then it would be a fourth radioactive primordial, though practically speaking it would still be more convenient to produce synthetically. Bismuth's half-life is so long that it is often classed with the 80 stable elements instead, since its radioactivity is not a cause for concern. The number of elements is smaller than the number of nuclides, because many of the primordial elements are represented by multiple isotopes. See chemical element for more information.
Naturally occurring stable nuclides
As noted, these number about 251. For a list, see the article list of elements by stability of isotopes. For a complete list noting which of the "stable" 251 nuclides may be in some respect unstable, see list of nuclides and stable nuclide. These questions do not impact the question of whether a nuclide is primordial, since all "nearly stable" nuclides, with half-lives longer than the age of the universe, are also primordial.
Radioactive primordial nuclides
Though it is estimated that about 35 primordial nuclides are radioactive (list below), it becomes very hard to determine the exact total number of radioactive primordials, because the total number of stable nuclides is uncertain. There are many extremely long-lived nuclides whose half-lives are still unknown; in fact, all nuclides heavier than dysprosium-164 are theoretically radioactive. For example, it is predicted theoretically that all isotopes of tungsten, including those indicated by even the most modern empirical methods to be stable, must be radioactive and can alpha decay, but this could only be measured experimentally for W. Likewise, all four primordial isotopes of lead are expected to decay to mercury, but the predicted half-lives are so long (some exceeding 10 years) that such decays could hardly be observed in the near future. Nevertheless, the number of nuclides with half-lives so long that they cannot be measured with present instruments—and are considered from this viewpoint to be stable nuclides—is limited. Even when a "stable" nuclide is found to be radioactive, it merely moves from the stable to the unstable list of primordials, and the total number of primordial nuclides remains unchanged. For practical purposes, these nuclides may be considered stable for all purposes outside specialized research.
List of 35 radioactive primordial nuclides and measured half-lives
These 35 primordial radionuclides are isotopes of 28 elements (cadmium, neodymium, osmium, samarium, tellurium, uranium, and xenon each have two primordial radioisotopes). These nuclides are listed in order of decreasing stability. Many of them are so nearly stable that they compete for abundance with stable isotopes of their respective elements.
For three elements (indium, tellurium, and rhenium) a very long-lived radioactive primordial nuclide is more abundant than a stable nuclide.
The longest-lived radionuclide known, Te, has a half-life of : 1.6 × 10 times the age of the Universe. Only four of these 35 nuclides have half-lives shorter than, or equal to, the age of the universe. Most of the other 30 have half-lives much longer. The shortest-lived primordial, U, has a half-life of 703.8 million years, about 1/6 the age of the Earth and Solar System. Many of these nuclides decay by double beta decay, though some like Bi decay by other means such as alpha decay.
At the end of the list, are two more nuclides: Sm and Pu. They have not been confirmed as primordial, but their half-lives are long enough that minute quantities should persist today.
List legends
| Physical sciences | Geochemistry | Earth science |
5292611 | https://en.wikipedia.org/wiki/Lithium%20oxide | Lithium oxide | Lithium oxide (O) or lithia is an inorganic chemical compound. It is a white solid. Although not specifically important, many materials are assessed on the basis of their Li2O content. For example, the Li2O content of the principal lithium mineral spodumene (LiAlSi2O6) is 8.03%.
Production
Lithium oxide forms along with small amounts of lithium peroxide when lithium metal is burned in the air and combines with oxygen at temperatures above 100 °C:
4Li + → 2.
Pure can be produced by the thermal decomposition of lithium peroxide, , at 450 °C
2 → 2 +
Structure
Solid lithium oxide adopts an antifluorite structure with four-coordinated Li+ centers and eight-coordinated oxides.
The ground state gas phase molecule is linear with a bond length consistent with strong ionic bonding. VSEPR theory would predict a bent shape similar to .
Uses
Lithium oxide is used as a flux in ceramic glazes; and creates blues with copper and pinks with cobalt. Lithium oxide reacts with water and steam, forming lithium hydroxide and should be isolated from them.
Its usage is also being investigated for non-destructive emission spectroscopy evaluation and degradation monitoring within thermal barrier coating systems. It can be added as a co-dopant with yttria in the zirconia ceramic top coat, without a large decrease in expected service life of the coating. At high heat, lithium oxide emits a very detectable spectral pattern, which increases in intensity along with degradation of the coating. Implementation would allow in situ monitoring of such systems, enabling an efficient means to predict lifetime until failure or necessary maintenance.
Lithium metal might be obtained from lithium oxide by electrolysis, releasing oxygen as by-product.
Reactions
Lithium oxide absorbs carbon dioxide forming lithium carbonate:
+ →
The oxide reacts slowly with water, forming lithium hydroxide:
+ → 2
| Physical sciences | Alkali oxide salts | Chemistry |
5293591 | https://en.wikipedia.org/wiki/Naver | Naver | Naver (; stylized as NAVER) is a South Korean online platform operated by the Naver Corporation. The company's products include a search engine, email hosting, blogs, maps, and mobile payment. It was founded in 1999 and was the first Korean company to develop their own search engine. Naver also created the world's first online Q&A platform, Knowledge iN.
As of September 2017, the search engine handled 74.7% of all web searches in South Korea and had 42 million registered users. More than 25 million Koreans have Naver as the start page on their default browser, and the mobile application has 28 million daily visitors.
History
Naver was the first Korean web provider to develop its own search engine. The company was founded on June 2, 1999 and is headquartered in Seongnam, South Korea. Naver is a combination of 'navigate', which means to navigate the Internet, and the suffix '-er', which means person, which means a person who navigates the vast ocean of information on the Internet. The Naver provides community services including blogs and cafes, other convenient services such as knowledge, shopping, maps, books, e-mail and naver tool bar.
In August 2000, Naver launched its 'comprehensive search' service, which allows users to get a variety of results from a single search query on one page, organized by type, including blogs, websites, images, and web communities. It was the world's first operator to introduce the comprehensive search feature, which compiles search results from various categories and presents them on a single page.
Naver became an early pioneer in user-generated content through the creation of the 'Knowledge iN' (네이버 지식인)' service in 2002. In Knowledge iN, users can pose questions on any subject and select from answers provided by other users, awarding points to those who give the best answers. Knowledge iN was launched three years before Yahoo! introduced its similar 'Yahoo! Answers' service and now possesses a database of over 200 million answers.
Bradley Horowitz, former Vice President of Product Strategy at Yahoo!, has cited Knowledge iN as the inspiration for Yahoo! Answers, which was launched three years after Naver introduced the original service.
Over the years, Naver has continued to expand its services. Naver Blog started with the name 'paper' in June 2003 and evolved to 'blog' in October 2003. 'Webtoon' is a South Korean webtoon publisher launched in 2004 by Naver.
In May 2005, Naver started the Real-time Search service. The rising search terms was a chart in Naver that highlighted trending topics in real-time. The top searches would appear on the portal main page reaching a large number of users.
From 2005 to 2007, Naver expanded its multimedia search services, including music and video search, as well as mobile search. Naver had launched a messaging app called Line for the South Korean market in February 2011.
In 2019, Naver reorganized its mobile version of the main screen, excluding search windows and some menus. In response, more than 3,000 comments opposing the change were posted.
In October 2023, Naver announced the beta release of its AI chatbot service, 'CLOVA X' (클로바X).
Naver mobile map application has begun providing real-time service on natural disaster information starting this September 15, 2024.
Services
Naver Dictionary
Naver Dictionary was launched in 1999, alongside Naver. It initially only supported Korean and English. , it supports 67 languages. The dictionary aggregates results from a number of other dictionaries, including Urimalsaem, which is operated by the National Institute of Korean Language. It also aggregates results from English-language dictionaries such as Oxford Dictionary of English and the Collins English Dictionary. It also operates an open-source dictionary called Open Dictionary PRO (ODP).
Junior Naver
Junior Naver (쥬니어 네이버), also known as Juniver (쥬니버), is a portal website for children, similar to Yahooligans. Junior Naver offers services such as avatars, educational content, quizzes, videos, Q&A, and a homework helper. It uses a panel of experts and educators to filter out harmful content, ensuring a safe internet environment for children. With its competitors Daum Kids and Yahoo Kids having closed down, Junior Naver is now the only children's portal site operating in Korea.
Knowledge iN
Knowledge iN (지식iN), formerly Knowledge Search (지식검색), is an online Q&A platform launched in October 2002. The tool allows users to ask any question and receive answers from other users. Naver became an early pioneer in user-generated content through the creation of the 'Knowledge iN' (네이버 지식인)' service in 2002. Bradley Horowitz, former Vice President of Product Strategy at Yahoo!, has cited Knowledge iN as the inspiration for Yahoo! Answers, which was launched three years after Naver introduced the original service.
Naver encourages unauthorized publishing to attract users to its Knowledge iN service. This contributes to poor-quality content on Knowledge iN, as previous answers to questions remain unchanged and old questions can only be minimally modified by other users. Criticism is also growing due to Naver's unilateral control over comments and its editing of information that may contradict its political positions. The 10s and 20s accounted for 70% of the questioners and 50% of the respondents. In April 2024, the service was reorganized so that answers could be continuously registered without closing questions, and new answers could be registered even if there were adopted answers. In addition, the limit on the number of additional questions and additional answers has disappeared. Even if you are not a questioner or an answerer, you can express your support for the answers you like.
Naver Webtoon
Naver Webtoon (네이버 웹툰), later simply WEBTOON, is a webcomic platform where users have free access to a variety of webtoons created by professional artists. Users can also pay publishers to view comic books and genre fiction content online. Naver has incorporated a 'Challenge' section that allows amateurs to post and promote their works.
Naver Cafe
Naver Cafe (네이버 카페) is a service that allows Naver users to create their own internet communities. As of May 2017, 10.5 million cafes were active. Each person can create up to 300 cafes.
From January to August 2024, Naver Cafe's MAU averaged 30 million, up about 10% from the previous year.
Naver Blog
Naver Blog (네이버 블로그) started with the name 'paper' in June 2003 and evolved to 'blog' in October 2003. It had 23 million users as of April 2016. In 2023, there were 1.26 million new users.
For two weeks from May 1, 2021, Naver held an event to pay up to 16,000 won to people who posted on Naver's blog every day. However, this event ended early due to several incidents involving people with multiple IDs. Many participants in the event criticized Naver's response.Naver announced on its official blog that it would resume its "Today's Diary Challenge (#오늘일기챌린지)" event, which ended early in three days, from May 24. However, only those who participated in the previously discontinued event (who completed the three-day record) can participate in the event.On May 13, 2021, Naver announced that it would display profile pictures along with comments posted on Naver news articles. Previously, only the first four digits of the author's ID were disclosed, making it difficult to identify users. Naver implemented this change expecting it to facilitate user recognition and address issues with malicious comments. Critics, however, criticized the company for censoring comments.
Naver NOW
Naver NOW (formerly Naver TV) is a video streaming and sharing platform that primarily offers web dramas distributed by Naver. Naver NOW replaced the Naver TV mobile app, while Naver TV continues to serve as a web portal.
Naver Pay
In June 2015, Naver launched its own payment service, Naver Pay, which allows mobile payment service and online checkout. It is Naver's second mobile payment service after Line Pay. Naver Pay is the most widely used mobile payment service in South Korea.
Naver Mail
Naver Mail (네이버 메일) is an email service available to all Naver users. Each user is provided with up to 5GB of storage.
Naver Shopping Live
Naver Shopping Live (네이버 쇼핑 라이브) is a live commerce platform operated by Naver. Broadcasts can be transmitted via mobile phones or camcorders, with a resolution of 1080*1920 and a vertical format similar to YouTube shorts. It is currently the platform with the highest number of consumers in Korea. As of 2023, 73.6% of mobile shoppers in South Korea use Naver Shopping Live as their primary platform.
PRISM Live Studio
PRISM Live Studio (프리즘 라이브 스튜디오) is a live streaming application available for both mobile and PC users. Streamers can simultaneously broadcast to multiple platforms, a practice known as simulcasting, with support for up to 1080p HD resolution without increasing network usage. Supported platforms include YouTube, Facebook, Twitch, Periscope, V Live, Naver TV, afreecaTV, KakaoTV, and RTMP channels. The application can also be utilized for video editing purposes.
Naver Papago
In July 2017, Naver launched Papago, which is an AI-based mobile translator that uses a large neural network technology named N2MT (Naver Neural Machine Translation). It can translate text and phrases in 15 different languages by analyzing context instead of statistical analysis.
Papago app has so far garnered over 16 million downloads.
Search engine features and restrictions
Naver has been criticized for abusing Real-time Search Terms (급상승 검색어) to manipulate public opinion.. As a result, Naver abolished Real-time Search Terms on February 25, 2021.
| Technology | Portals/Platform sites | null |
5298371 | https://en.wikipedia.org/wiki/%CE%9418O | Δ18O | {{DISPLAYTITLE:δ18O}}
In geochemistry, paleoclimatology and paleoceanography δ18O or delta-O-18 is a measure of the deviation in ratio of stable isotopes oxygen-18 (18O) and oxygen-16 (16O).
It is commonly used as a measure of the temperature of precipitation, as a measure of groundwater/mineral interactions, and as an indicator of processes that show isotopic fractionation, like methanogenesis.
In paleosciences, 18O:16O data from corals, foraminifera and ice cores are used as a proxy for temperature.
It is defined as the deviation in "per mil" (‰, parts per thousand) between a sample and a standard:
‰
where the standard has a known isotopic composition, such as Vienna Standard Mean Ocean Water (VSMOW). The fractionation can arise from kinetic, equilibrium, or mass-independent fractionation.
Mechanism
Foraminifera shells are composed of calcium carbonate (CaCO3) and are found in many common geological environments. The ratio of 18O to 16O in the shell is used to indirectly determine the temperature of the surrounding water at the time the shell was formed. The ratio varies slightly depending on the temperature of the surrounding water, as well as other factors such as the water's salinity, and the volume of water locked up in ice sheets.
also reflects local evaporation and freshwater input, as rainwater is 16O-enriched—a result of the preferential evaporation of the lighter 16O from seawater. Consequently, the surface ocean contains greater proportions of 18O around the subtropics and tropics where there is more evaporation, and lesser proportions of 18O in the mid-latitudes where it rains more.
Similarly, when water vapor condenses, heavier water molecules holding 18O atoms tend to condense and precipitate first. The water vapor gradient heading from the tropics to the poles gradually becomes more and more depleted of 18O. Snow falling in Canada has much less H218O than rain in Florida; similarly, snow falling in the center of ice sheets has a lighter signature than that at its margins, since heavier 18O precipitates first.
Changes in climate that alter global patterns of evaporation and precipitation therefore change the background ratio.
Solid samples (organic and inorganic) for oxygen isotope analysis are usually stored in silver cups and measured with pyrolysis and mass spectrometry. Researchers need to avoid improper or prolonged storage of the samples for accurate measurements.
Extrapolation of temperature
Based on the simplifying assumption that the signal can be attributed to temperature change alone, with the effects of salinity and ice volume change ignored, Epstein et al. (1953) estimated that
a increase of 0.22‰ is equivalent to a cooling of 1 °C (or 1.8 °F).
More precisely, Epstein et al. (1953) give a quadratic extrapolation for the temperature, as
where T is the temperature in °C (based on a least-squares fit for a range of temperature values between 9 °C and 29 °C, with a standard deviation of ±0.6 °C, and δ is δ18O for a calcium carbonate sample).
Paleoclimatology
Ice cores
δ18O can be used with ice cores to determine the temperature from when the ice was formed.
Lisiecki and Raymo (2005) used measurements of δ18O in benthic foraminifera from 57 globally distributed deep sea sediment cores, taken as a proxy for the total global mass of glacial ice sheets, to reconstruct the climate for the past five million years.
The stacked record of the 57 cores was orbitally tuned to an orbitally driven ice model, the Milankovitch cycles of 41 ky (obliquity), 26 ky (precession) and 100 ky (eccentricity), which are all assumed to cause orbital forcing of global ice volume. Over the past million years, there have been a number of very strong glacial maxima and minima, spaced by roughly 100 ky.
As the observed isotope variations are similar in shape to the temperature variations recorded for the past 420 ky at Vostok Station, the figure shown on the right aligns the values of δ18O (right scale) with the reported temperature variations from the Vostok ice core (left scale), following Petit et al. (1999).
Biomineralized tissues
δ18O from biomineralized tissues may also be used in reconstructing past environmental conditions. In vertebrates, apatite from bone mineral, tooth enamel and dentin contains phosphate [PO4]3− groups which may preserve the oxygen isotope ratios of environmental water. Fractionation of oxygen isotopes in these tissues may be affected by biological factors such as body temperature and diet.
| Physical sciences | Geochemistry | Earth science |
5298400 | https://en.wikipedia.org/wiki/Arandaspis | Arandaspis | Arandaspis prionotolepis is an extinct species of jawless fish that lived in the Ordovician period, about 480 to 470 million years ago. Its remains were found in the Stairway Sandstone near Alice Springs, Australia in 1959, but it was not determined that they were the oldest known vertebrates until the late 1960s. Arandaspis is named after a local Indigenous Australian people, the Aranda (now currently called Arrernte).
Description
Arandaspis is estimated to reach around long, with a body covered in rows of knobbly armoured scutes. The front of the body and the head were protected by hard plates with openings for the eyes, nostrils and gills. It probably was a filter-feeder. Morphology of trunk and tail is unknown. According to comparison with other early ostracoderms, it would lacked paired fins and caudal fin would be simple shape, although another arandaspid Sacabambaspis had the tail that consist dorsal and ventral webs and an elongated notochordal lobe.
| Biology and health sciences | Prehistoric agnathae and early chordates | Animals |
1455717 | https://en.wikipedia.org/wiki/A%C3%A7a%C3%AD%20palm | Açaí palm | The açaí palm (, , from Nheengatu asai), Euterpe oleracea, is a species of palm tree (Arecaceae) cultivated for its fruit (açaí berries, or simply açaí), hearts of palm (a vegetable), leaves, and trunk wood. Global demand for the fruit has expanded rapidly in the 21st century, and the tree is cultivated for that purpose primarily.
The species is native to eastern Amazonia, especially in Brazil, mainly in swamps and floodplains. Açaí palms are tall, slender trees growing to more than tall, with pinnate leaves up to long. The fruit is small, round, and black-purple in color. The fruit became a staple food in floodplain areas around the 18th century, but its consumption in urban areas and promotion as a health food only began in the mid 1990s along with the popularization of other Amazonian fruits outside the region.
Name
The common name comes from the Portuguese adaptation of the Tupian word , meaning "[fruit that] cries or expels water". The importance of the fruit as a staple food in the Amazon River delta gives rise to the local legend of how the plant got its name. The folklore says that chief Itaqui ordered all newborns put to death owing to a period of famine. When his own daughter gave birth and the child was sacrificed, she cried and died beneath a newly sprouted tree. The tree fed the tribe and was called açaí because that was the daughter's name (Iaçá) spelled backwards.
Its specific epithet oleracea means "vegetable" in Latin and is a form of ().
Fruit
The fruit, commonly known as açaí or açaí berry, is a small, round, black-purple drupe about in circumference, similar in appearance to a grape, but smaller and with less pulp and produced in branched panicles of 500 to 900 fruits. The exocarp of the ripe fruits is a deep purple color, or green, depending on the kind of açaí and its maturity. The mesocarp is pulpy and thin, with a consistent thickness of or less. It surrounds the voluminous and hard endocarp, which contains a single large seed about in diameter. The seed makes up about 60–80% of the fruit. The palm bears fruit year round but the berry cannot be harvested during the rainy season.
Cultivation
There are two harvests: one is normally between January and June, while the other is between August and December, producing larger volumes. In 2022, the state of Pará, which accounts for 90% of Brazil's total açaí economy, produced of açaí berries, generating US$26 million in revenue. The 2022 production was 209 times greater than the volume produced in 2012.
Child labor concern
Children as young as 13 years old are employed as laborers to harvest the fruit, using machetes to clear paths in the rainforest, and climbing trees up to tall without harnesses to collect berries in the canopy, a process leading to falls and severe injuries in some children.
Cultivars
Few named cultivars exist, and varieties differ mostly in the nature of the fruit:
Branco ("White") is a rare variety local to the Amazon estuary in which the berries do not change color, but remain green when ripe. This is believed to be due to a recessive gene since only about 30% of 'Branco' palm seeds mature to express this trait.
BRS-Pará was developed in 2004 by the Brazilian Agricultural Research Agency. The pulp yield ranges from 15% to 25%.
BRS Pai d'Égua is the newest cultivar developed by the Brazilian Agricultural Research Agency.
Anthocyanins
Anthocyanins define the blue pigmentation of açaí and the antioxidant capacity of the plant's natural defense mechanisms and in laboratory experiments in vitro. Anthocyanins in açaí accounted for only about 10% of the overall antioxidant capacity in vitro. The Linus Pauling Institute and European Food Safety Authority state that "the relative contribution of dietary flavonoids to (...) antioxidant function in vivo is likely to be very small or negligible". Unlike in controlled test tube conditions, anthocyanins have been shown to be poorly conserved (less than 5%) in vivo, and most of what is absorbed exists as chemically modified metabolites destined for rapid excretion.
A powdered preparation of freeze-dried açaí fruit pulp and skin was shown to contain cyanidin 3-O-glucoside and cyanidin 3-O-rutinoside as major anthocyanins (3.19 mg/g). The powdered preparation was also reported to contain twelve flavonoid-like compounds, including homoorientin, orientin, taxifolin deoxyhexose, isovitexin, scoparin, as well as proanthocyanidins (12.89 mg/g), and low levels of resveratrol (1.1 μg/g).
Nutritional content
A powdered preparation of freeze-dried açaí fruit pulp and skin was reported to contain (per 100 g of dry powder) 534 calories, 52 g carbohydrates, 8 g protein, and 33 g total fat. The carbohydrate portion included 44 g of dietary fiber with low sugar levels, and the fat portion consisted of oleic acid (56% of total fats), palmitic acid (24%), and linoleic acid (13%). The powder was also shown to contain (per 100 g) negligible vitamin C, 260 mg calcium, 4 mg iron, and 1002 IU vitamin A.
Marketing
In the 1980s, the Brazilian Gracie family marketed açaí as an energy drink or as crushed fruit served with granola and bananas; this demand led to the building of cottage industries and processing plants to pulp and freeze açaí for export.
Scams
In the early 2000s, numerous companies advertised açaí products online, with many ads featuring counterfeit testimonials and products. In 2009, açaí scams were ranked No. 1 on the U.S. Federal Trade Commission's "scams and rip-offs" list, so that by 2011 sales of açaí flattened as the fad waned.
According to the Washington, D.C.–based Center for Science in the Public Interest thousands of consumers had trouble stopping recurrent charges on their credit cards when they canceled free trials of some açai-based products. In 2003, American celebrity doctor Nicholas Perricone included açaí berries among "superfoods", but such extravagant marketing claims regarding açaí as miracle cures for everything from obesity to attention-deficit disorder were challenged in subsequent studies.
The FTC handed down an $80 million judgement in January 2012 against five companies that were marketing açaí berry supplements with fraudulent claims that their products promoted weight loss and prevented colon cancer. One company, Central Coast Nutraceuticals, was ordered to pay a $1.5 million settlement.
Production
Brazil is a major producer, particularly in the state of Pará, which alone in 2019 produced more than 1.2 million tons of açaí, an amount equal to 95% of Brazil's total.
Uses
As a food product
Fresh açaí has been consumed as a dietary staple in the region around the Amazon river delta for centuries. The fruit is processed into pulp for supply to food product manufacturers or retailers, sold as frozen pulp, juice, or an ingredient in various products from beverages, including grain alcohol, smoothies, foods, cosmetics and supplements. In Brazil, it is commonly eaten as .
In a study of three traditional Caboclo populations in the Brazilian Amazon, açaí palm was described as the most important plant species because the fruit makes up a major component of their diet, up to 42% of the total food intake by weight.
Açaí na tigela (known in English as açaí bowl) is a Brazilian dessert made from frozen açaí berry purée, served in a bowl and topped with other fruit and granola.
Dietary supplement
As of 2008, no açaí products have been evaluated by the FDA, and their efficacy is doubtful.
As of 2009, there is no scientific evidence that açaí consumption affects body weight, promotes weight loss or has any positive health effect.
Açaí oil
Açaí oil is suitable for cooking or as a salad dressing, but is mainly used in cosmetics as shampoos, soaps or skin moisturizers.
The oil compartments in açaí fruit contain polyphenols such as procyanidin oligomers and vanillic acid, syringic acid, p-hydroxybenzoic acid, protocatechuic acid, and ferulic acid, which were shown to degrade substantially during storage or exposure to heat. Although these compounds are under study for potential health effects, there remains no substantial evidence that açaí polyphenols have any effect in humans. Açaí oil is green in color, has a bland aroma, and is high in oleic and palmitic fatty acids.
Other uses
Leaves of the palm may be made into hats, mats, baskets, brooms and roof thatch for homes, and trunk wood, resistant to pests, for building construction. Tree trunks may be processed to yield dietary minerals.
Comprising 80% of the fruit mass, açaí seeds may be ground for livestock food or as a component of organic soil for plants. Planted seeds are used for new palm tree stock, which, under the right growing conditions, can require months to form seedlings. Seeds may become waste in landfills or used as fuel for producing bricks.
Research
Orally administered açaí has been tested as a contrast agent for magnetic resonance imaging (MRI) of the gastrointestinal system. Its anthocyanins have also been characterized for stability as a natural food coloring agent.
Gallery
| Biology and health sciences | Other culinary fruits | Plants |
1455879 | https://en.wikipedia.org/wiki/Tusk%20shell | Tusk shell | Scaphopoda (plural scaphopods , from Ancient Greek σκᾰ́φης skáphē "boat" and πούς poús "foot"), whose members are also known as tusk shells or tooth shells, are a class of shelled marine invertebrates belonging to the phylum Mollusca with worldwide distribution and are the only class of exclusively infaunal marine molluscs. Shells of species within this class range in length (with Fissidentalium metivieri as the longest). Members of the order Dentaliida tend to be larger than those of the order Gadilida.
These molluscs live in soft substrates offshore (usually not intertidally). Because of this subtidal habitat and the small size of most species, many beachcombers are unfamiliar with them; their shells are not as common or as easily visible in the beach drift as the shells of sea snails and clams.
Molecular data suggest that the scaphopods are a sister group to the cephalopods, although higher-level molluscan phylogeny remains unresolved.
Classification
The group is composed of two subtaxa, the Dentaliida (which may be paraphyletic) and the monophyletic Gadilida.
The differences between the two orders is subtle and hinges on size and on details of the radula, shell, and foot. Specifically, the Dentaliids are the physically larger of the two families, and possess a shell that tapers uniformly from anterior (widest) to posterior (narrowest); they also have a foot which consists of one central and two lateral lobes and which bends into the shell when retracted. The Gadilids, on the other hand, are much smaller, have a shell whose widest portion is slightly posterior to its aperture, and have a foot which is disk-like and fringed with tentacles which inverts into itself when retracted (in this state resembling a pucker rather than a disk).
According to the World Register of Marine Species:
Dentaliida da Costa, 1776
family Anulidentaliidae Chistikov, 1975 – 3 genera
family Calliodentaliidae – 1 genus
family Dentaliidae Children, 1834 – 14 genera
family Fustiariidae Steiner, 1991 – 1 genus
family Gadilinidae Chistikov, 1975 – 3 genera
family Laevidentaliidae Palmer, 1974 – 1 genus
family Omniglyptidae Chistikov, 1975 – 1 genus
family Rhabdidae Chistikov, 1975 – 1 genus
Gadilida Starobogatov, 1974
sub-order Entalimorpha Steiner, 1992
family Entalinidae Chistikov, 1979 – 9 genera
sub-order Gadilimorpha Steiner, 1992
family Gadilidae Stoliczka, 1868 – 8 genera
family Pulsellidae Scarabino in Boss, 1982 – 3 genera
family Wemersoniellidae Scarabino, 1986 – 2 genera
Evolution
Fossil record
There is a good fossil record of scaphopods from the Mississippian onwards, making them the youngest molluscan class.
The Ordovician Rhytiodentalium kentuckyensis has been interpreted as an early antecedent of the scaphopods, implying an evolutionary succession from ribeirioid rostroconch molluscs such as Pinnocaris. However, a competing hypothesis suggests a Devonian/Carboniferous origin from a non-mineralized ancestor, or from a more derived, Devonian, conocardioid rostroconch.
Phylogeny
The scaphopods are largely agreed to be members of the Conchifera, however their phylogenetic relationship with the other members of this subphylum remains contentious. The Diasoma concept proposes a clade of scaphopods and bivalves based on their shared infaunal lifestyle, burrowing foot, and possession of a mantle and shell. Pojeta and Runnegar proposed the extinct Rostroconchia as the stem group of the Diasoma. An alternative hypothesis proposes the cephalopods and gastropods as sister to the scaphopods with helcionellids as the stem group. A review of deep molluscan phylogeny in 2014 found more support for the scaphopods, gastropods, or cephalopods than for scaphopods and bivalves, thus the shared body features of scaphopods and bivalves may be convergent adaptations due to similar lifestyles. Analysis of the scaphopod nervous system demonstrated that both scaphopods and cephalopods share a similar nervous system structure, with ventrally shifted pedal nerves and lateral nerves that extend dorsally. These similarities led to the conclusion that scaphopods are sister to the cephalopods with gastropods as sister to them both. More recent research, including the sequenced genome of tusk shells, support the Diasoma model with bivalves as the sister group.
Orientation
The morphological shape of the scaphopod body makes it difficult to orient it satisfactorily. As a result, researchers have often disagreed as to which direction is anterior/ posterior and which is ventral/ dorsal. According to Shimek and Steiner, "[t]he apex of the shell and mantle are anatomically dorsal, and the large aperture is ventral and anterior. Consequently, the concave side of the shell and viscera are anatomically dorsal. The convex side has to be divided into anteriorly ventral and dorsally posterior portions, with the anus as the demarcation. Functionally, as in cephalopods, the large aperture with the foot is anterior, the apical area posterior, the concave side dorsal and the convex side ventral."
Anatomy
Shells
The shells of the members of the Gadilida are usually glassy-smooth and narrow, with a reduced aperture. This along with other structures of their anatomy allows them to move with surprising speed through loose sediment to escape potential bottom-dwelling predators.
The Dentalids, by contrast, tend to have strongly ribbed and rough shells. When they sense vibrations anywhere around them, their defensive response is to freeze. This makes them harder to detect by animals such as ratfish, which can sense the electrical signals given off by the most minute muscle movement.
Mantle
The mantle of a scaphopod is entirely within the shell. The foot extends from the larger end of the shell, and is used to burrow through the substrate. The scaphopod positions itself head down in the substrate, with the apical end of the shell (at the rear of the animal's body) projecting upward. This end seldom appears above the level of the substrate, however, as doing so exposes the animal to numerous predators. Most adult scaphopods live their lives entirely buried within the substrate.
Water enters the mantle cavity through the apical aperture, and is wafted along the body surface by cilia. There are no gills; the entire surface of the mantle cavity absorbs oxygen from the water. Unlike most other molluscs, there is no continuous flow of water with a separate exhalant stream. Instead, deoxygenated water is expelled rapidly back through the apical aperture through muscular action once every ten to twelve minutes.
Feeding and digestion
A number of minute tentacles around the foot, called captacula, sift through the sediment and latch onto bits of food, which they then convey to the mouth. The mouth has a grinding radula that breaks the bit into smaller pieces for digestion. The radulae and cartilaginous oral bolsters of the Gadilidae are structured like zippers where the teeth actively crush the prey by opening and closing on it repeatedly, while the radulae and bolsters of the Dentaliidae work rachet-like to pull the prey into the esophagus, sometimes whole.
The massive radula of the scaphopods is the largest such organ relative to body size of any mollusc (among whom, except for the bivalves, the presence of which is a defining characteristic). The remainder of the digestive system consists of a digestive diverticulum, esophagus, stomach, and intestine. A digestive gland secretes enzymes into the stomach, but, unlike some other molluscs, does not digest the food directly itself. The anus opens on the ventral/ underside of the animal, roughly in the middle of the mantle cavity.
Vascular system
The scaphopod vascular system is rudimentary lacking heart auricles as well as corresponding ctenidia (gills) and blood vessels; the blood is held in sinuses throughout the body cavity, and is pumped through the body by the rhythmic action of the foot. The heart, a characteristic feature of all other groups of mollusca, has been considered totally lost or reduced to a thin fold of the pericardium; however, according to more recent studies, the muscular, regularly beating perianal blood sinus is homologous to the ventricle and is therefore considered the scaphopod heart.
Metabolic waste is excreted through a pair of nephridia close to the anus. The tusk shells appear to be the only extant molluscs which completely lack the otherwise standard molluscan reno-pericardial apertures. Furthermore, they also appear to be the only molluscs with openings that directly connect the hemocoel with the surrounding water (through two "water pores" located near the nephridial openings). These openings may serve to allow the animal to relieve internal pressure by ejecting body fluid (blood) during moments of extreme muscular contraction of the foot.
Nervous system
The nervous system is generally similar to that of cephalopods. One pair each of cerebral and pleural ganglia lie close to the oesophagus, and effectively form the animal's brain.
A separate set of pedal ganglia lie in the foot, and a pair of visceral ganglia are set further back in the body, and connect to pavilion ganglia via long connectives. Radular and sub-radular ganglia are also present, as are statocysts with staticonia. Scaphopods have no eyes, no osphradia, or other distinct sensory organs. However, scaphopods do possess genes involved in photoreceptor formation and function implying scaphopods may have had eyes that degenerated over evolutionary time.
Reproduction and development
Scaphopods have separate sexes, and external fertilisation. They have a single gonad occupying much of the posterior part of the body, and shed their gametes into the water through the nephridium.
Once fertilized, the eggs hatch into a free-living trochophore larva, which develops into a veliger larva that more closely resembles the adult, but lacks the extreme elongation of the adult body. The three-lobed foot originates prior to metamorphosis while the cephalic tentacles develop post metamorphosis. Scaphopods remain univalved throughout their morphogenesis contrary to bivalves.
Ecology
Tusk shells live in seafloor sediment, feeding primarily on foraminiferans; some supplement this with vegetable matter.
Human use
The shells of Dentalium hexagonum and Dentalium pretiosum were strung on thread and used by the natives of the Pacific Northwest as shell money. Dentalium shells were also used to make belts and headdresses by the Natufian culture of the Middle East, and are a possible indicator of early social stratification.
| Biology and health sciences | Mollusks | Animals |
1456984 | https://en.wikipedia.org/wiki/Organic%20synthesis | Organic synthesis | Organic synthesis is a branch of chemical synthesis concerned with the construction of organic compounds. Organic compounds are molecules consisting of combinations of covalently-linked hydrogen, carbon, oxygen, and nitrogen atoms. Within the general subject of organic synthesis, there are many different types of synthetic routes that can be completed including total synthesis, stereoselective synthesis, automated synthesis, and many more. Additionally, in understanding organic synthesis it is necessary to be familiar with the methodology, techniques, and applications of the subject.
Total synthesis
A total synthesis refers to the complete chemical synthesis of molecules from simple, natural precursors. Total synthesis is accomplished either via a linear or convergent approach. In a linear synthesis—often adequate for simple structures—several steps are performed sequentially until the molecule is complete; the chemical compounds made in each step are called synthetic intermediates. Most often, each step in a synthesis is a separate reaction taking place to modify the starting materials. For more complex molecules, a convergent synthetic approach may be better suited. This type of reaction scheme involves the individual preparations of several key intermediates, which are then combined to form the desired product.
Robert Burns Woodward, who received the 1965 Nobel Prize for Chemistry for several total syntheses including his synthesis of strychnine, is regarded as the grandfather of modern organic synthesis. Some latter-day examples of syntheses include Wender's, Holton's, Nicolaou's, and Danishefsky's total syntheses of the anti-cancer drug paclitaxel (trade name Taxol).
Methodology and applications
Before beginning any organic synthesis, it is important to understand the chemical reactions, reagents, and conditions required in each step to guarantee successful product formation. When determining optimal reaction conditions for a given synthesis, the goal is to produce an adequate yield of pure product with as few steps as possible. When deciding conditions for a reaction, the literature can offer examples of previous reaction conditions that can be repeated, or a new synthetic route can be developed and tested. For practical, industrial applications additional reaction conditions must be considered to include the safety of both the researchers and the environment, as well as product purity.
Synthetic techniques
Organic Synthesis requires many steps to separate and purify products. Depending on the chemical state of the product to be isolated, different techniques are required. For liquid products, a very common separation technique is liquid–liquid extraction and for solid products, filtration (gravity or vacuum) can be used.
Liquid–liquid extraction
Liquid–liquid extraction uses the density and polarity of the product and solvents to perform a separation. Based on the concept of "like-dissolves-like", non-polar compounds are more soluble in non-polar solvents, and polar compounds are more soluble in polar solvents. By using this concept, the relative solubility of compounds can be exploited by adding immiscible solvents into the same flask and separating the product into the solvent with the most similar polarity. Solvent miscibility is of major importance as it allows for the formation of two layers in the flask, one layer containing the side reaction material and one containing the product. As a result of the differing densities of the layers, the product-containing layer can be isolated and the other layer can be removed.
Heated reactions and reflux condensers
Many reactions require heat to increase reaction speed. However, in many situations increased heat can cause the solvent to boil uncontrollably which negatively affects the reaction, and can potentially reduce product yield. To address this issue, reflux condensers can be fitted to reaction glassware. Reflux condensers are specially calibrated pieces of glassware that possess two inlets for water to run in and out through the glass against gravity. This flow of water cools any escaping substrate and condenses it back into the reaction flask to continue reacting and ensure that all product is contained. The use of reflux condensers is an important technique within organic syntheses and is utilized in reflux steps, as well as recrystallization steps.
When being used for refluxing a solution, reflux condensers are fitted and closely observed. Reflux occurs when condensation can be seen dripping back into the reaction flask from the reflux condenser; 1 drop every second or few seconds.
For recrystallization, the product-containing solution is equipped with a condenser and brought to reflux again. Reflux is complete when the product-containing solution is clear. Once clear, the reaction is taken off heat and allowed to cool which will cause the product to re-precipitate, yielding a purer product.
Gravity and vacuum filtration
Solid products can be separated from a reaction mixture using filtration techniques. To obtain solid products a vacuum filtration apparatus can be used.
Vacuum filtration uses suction to pull liquid through a Büchner funnel equipped with filter paper, which catches the desired solid product. This process removes any unwanted solution in the reaction mixture by pulling it into the filtration flask and leaving the desired product to collect on the filter paper.
Liquid products can also be separated from solids by using gravity filtration. In this separatory method, filter paper is folded into a funnel and placed on top of a reaction flask. The reaction mixture is then poured through the filter paper, at a rate such that the total volume of liquid in the funnel does not exceed the volume of the funnel. This method allows for the product to be separated from other reaction components by the force of gravity, instead of a vacuum.
Stereoselective synthesis
Most complex natural products are chiral, and the bioactivity of chiral molecules varies with the enantiomer. Some total syntheses target racemic mixtures, which are mixtures of both possible enantiomers. A single enantiomer can then be selected via enantiomeric resolution.
As chemistry has developed methods of stereoselective catalysis and kinetic resolution have been introduced whereby reactions can be directed, producing only one enantiomer rather than a racemic mixture. Early examples include stereoselective hydrogenations (e.g., as reported by William Knowles and Ryōji Noyori) and functional group modifications such as the asymmetric epoxidation by Barry Sharpless; for these advancements in stereochemical preference, these chemists were awarded the Nobel Prize in Chemistry in 2001. Such preferential stereochemical reactions give chemists a much more diverse choice of enantiomerically pure materials.
Using techniques developed by Robert B. Woodward paired with advancements in synthetic methodology, chemists have been able synthesize stereochemically selective complex molecules without racemization. Stereocontrol provides the target molecules to be synthesized as pure enantiomers (i.e., without need for resolution). Such techniques are referred to as stereoselective synthesis.
Synthesis design
Many synthetic procedures are developed from a retrosynthetic framework, a type of synthetic design developed by Elias James Corey, for which he won the Nobel Prize in Chemistry in 1990. In this approach, the synthesis is planned backwards from the product, obliging to standard chemical rules. Each step breaks down the parent structure into achievable components, which are shown via the use of graphical schemes with retrosynthetic arrows (drawn as ⇒, which in effect, means "is made from"). Retrosynthesis allows for the visualization of desired synthetic designs.
Automated organic synthesis
A recent development within organic synthesis is automated synthesis. To conduct organic synthesis without human involvement, researchers are adapting existing synthetic methods and techniques to create entirely automated synthetic processes using organic synthesis software. This type of synthesis is advantageous as synthetic automation can increase yield with continual "flowing" reactions. In flow chemistry, substrates are continually fed into the reaction to produce a higher yield. Previously, this type of reaction was reserved for large-scale industrial chemistry but has recently transitioned to bench-scale chemistry to improve the efficiency of reactions on a smaller scale.
Currently integrating automated synthesis into their work is SRI International, a nonprofit research institute. Recently SRI International has developed Autosyn, an automated multi-step chemical synthesizer that can synthesize many FDA-approved small molecule drugs. This synthesizer demonstrates the versatility of substrates and the capacity to potentially expand the type of research conducted on novel drug molecules without human intervention.
Automated chemistry and the automated synthesizers used demonstrate a potential direction for synthetic chemistry in the future.
Characterization
Necessary to organic synthesis is characterization. Characterization refers to the measurement of chemical and physical properties of a given compound, and comes in many forms. Examples of common characterization methods include: nuclear magnetic resonance (NMR), mass spectrometry, Fourier-transform infrared spectroscopy (FTIR), and melting point analysis. Each of these techniques allow for a chemist to obtain structural information about a newly synthesized organic compound. Depending on the nature of the product, the characterization method used can vary.
Relevance
Organic synthesis is an important chemical process that is integral to many scientific fields. Examples of fields beyond chemistry that require organic synthesis include the medical industry, pharmaceutical industry, and many more. Organic processes allow for the industrial-scale creation of pharmaceutical products. An example of such a synthesis is Ibuprofen. Ibuprofen can be synthesized from a series of reactions including: reduction, acidification, formation of a Grignard reagent, and carboxylation.
In the synthesis of Ibuprofen proposed by Kjonass et al., p-isobutylacetophenone, the starting material, is reduced with sodium borohydride (NaBH4) to form an alcohol functional group. The resulting intermediate is acidified with HCl to create a chlorine group. The chlorine group is then reacted with magnesium turnings to form a Grignard reagent. This Grignard is carboxylated and the resulting product is worked up to synthesize ibuprofen.
This synthetic route is just one of many medically and industrially relevant reactions that have been created, and continued to be used.
| Physical sciences | Synthetic strategies | Chemistry |
1457471 | https://en.wikipedia.org/wiki/Tokay%20gecko | Tokay gecko | The tokay gecko (Gekko gecko) is a nocturnal arboreal gecko in the genus Gekko, the true geckos. It is native to Asia and some Pacific Islands.
Etymology
The word "tokay" is an onomatopoeia of the sound made by males of this species. The common and scientific names, as well as the family name Gekkonidae and the generic term "gecko" come from this species, too, from ge'kok in Javanese, corresponding to tokek in Malay.
Subspecies
Two subspecies are currently recognized:
G. g. gecko (Linnaeus, 1758) occurs in tropical Asia from northeastern India to eastern Indonesia.
G. g. azhari (Mertens, 1955) is found only in Bangladesh.
Distribution and habitat
This species is found in northeast India, Bhutan, Nepal, and Bangladesh; throughout Southeast Asia, including Cambodia, Thailand, the Philippines, Malaysia, Vietnam and Indonesia; and toward western New Guinea. Its native habitat is rainforests, where it lives on trees and cliffs, and it frequently adapts to rural human habitations, roaming walls and ceilings at night in search of insect prey. This is an introduced species in some areas outside its native range. It is established in Florida in the United States, Martinique, the islands of Belize, and possibly Hawaii. Increasing urbanization is reducing its range.
Whether the species is native but very uncommon in Taiwan or whether the rare reports of individuals since the 1920s are based on repeated anthropogenic translocations that may or may not have resulted in established populations by now is unclear.
Physical characteristics
The tokay gecko is a large nocturnal gecko, reaching a total length (including tail) of on average, but some grow as large as long. It is believed to be the third-largest species of gecko, after the giant leaf-tail gecko (Uroplatus giganteus) and New Caledonian giant gecko (Rhacodactylus leachianus). It is cylindrical, but somewhat flattened in body shape. The eyes have vertical pupils.
The skin is soft to the touch and is generally blue-gray with red or orange spots and speckles, but the animal can change the color of its skin to blend into the environment. The species is sexually dimorphic, with the males being more brightly colored and slightly larger than females.
They have two visual pigments: a "green" with λmax at 521 nm and a "blue" at 467 nm.
It is a strong climber with foot pads that can support the entire weight of its body on a vertical surface for a long period of time. Compared to other gecko species, the tokay gecko has a robust build, with a semiprehensile tail, a large head, and muscular jaws.
Behaviour
Tokay geckos are generally aggressive and territorial, and can inflict a strong bite. Though common in the pet trade, the strong bite of the tokay gecko makes it ill-suited for inexperienced keepers. In addition, the strength of the bite depends on the gecko's size; larger (usually male) tokay geckos are capable of piercing skin, which often results in immediate bleeding.
Females lay clutches of one or two hard-shelled eggs and guard them until they hatch. Communal nesting has been observed in Cambodia with as many as four adults guarding the clutch of about eleven eggs.
Diet
The tokay gecko feeds on insects and small vertebrates. In captivity, they usually feed on springtails, mealworms, cockroaches, crickets, grasshoppers, locusts, and pink mice. In a study conducted in Thailand, researchers noticed little variation in the diets of males, females, and juveniles, which was likely due to low insect availability in this area. In Southern Florida, roaches, caterpillars, spiders, and beetles were the most numerous taxa present and common to the greatest number of stomachs.
Call
The male's mating call, a loud croak, is variously described as sounding like token, gekk-gekk, tuck-too, túc-key, tou-kay or tokay. Most of the time, the call is often preceded by a quick "cackling", similar to the chirping sounds made by house geckos albeit much lower in pitch. When threatened or alarmed, tokay geckos usually "bark" while opening their mouth in a defensive posture.
The tokay gecko's call is also responsible for the name given to it by Filipino residents: "Tuko," and by U.S. soldiers during the Vietnam War, the "Fuck-you lizard".
Light and temperature can affect its vocalizations. The most frequent calling occurs in May at dusk, and the second peak of call frequency occurs in May at dawn. Vocalizations and associated behavior were strongly affected by ambient temperature in both the lab and field and could thus play a role in regulating animal energetic metabolism
Conservation and relationship with humans
The tokay gecko is culturally significant in many East Asian countries. Regional folklore has attributed supernatural powers to the gecko. In Southeast Asia it is a symbol of good luck and fertility. It is believed to be descended from dragons.
This species is poached for the medicinal trades in parts of Asia. The tokay gecko is an ingredient in traditional Chinese medicine, known as Ge Jie (蛤蚧). It is believed to nourish the kidneys and lungs, beliefs that are not substantiated by medical science. The animal remains highly sought after in China, Hong Kong, Taiwan, Vietnam, Malaysia, Singapore and other parts of Asia with Chinese communities, to the point where unscrupulous merchants have taken to disfiguring monitor lizards with prosthetics to pass them off as colossal tokay gecko specimens.
From 2009 to 2011, the poaching of tokay geckos intensified because of a short-lived belief that they were an effective HIV cure.
The tokay gecko is quickly becoming a threatened species in the Philippines because of indiscriminate hunting. Collecting, transporting and trading in geckos without a license can be punishable by up to 12 years in jail and a fine of up to ₱1 million under Republic Act 9147, in addition to other applicable international laws. However, the trade runs unchecked because of the sheer number of illegal traders and reports of lucrative deals. Chinese buyers and other foreign nationals are rumored to pay thousands of dollars for large specimens, because of their alleged medicinal value or as commodities in the illegal wildlife trade.
The species is protected under Appendix II of the Convention on International Trade in Endangered Species (CITES) meaning international trade (including in parts and derivatives) is subject to the CITES permitting system.
Captivity
Tokay geckos are becoming more popular as pets because of their striking colors and large size. Most of them are wild-caught imports, but captive-bred ones are becoming more common. Wild-caught adults can be difficult to keep because of their aggressive nature and powerful bite, but captive-bred juveniles can be less aggressive if handled from a young age.
When well cared for, tokay geckos can live up to 15–20 years.
| Biology and health sciences | Lizards and other Squamata | Animals |
1457625 | https://en.wikipedia.org/wiki/Caspian%20tiger | Caspian tiger | The Caspian tiger was a Panthera tigris tigris population native to eastern Turkey, northern Iran, Mesopotamia, the Caucasus around the Caspian Sea, Central Asia to northern Afghanistan and the Xinjiang region in western China. Until the Middle Ages, it was also present in southern Russia. It inhabited sparse forests and riverine corridors in this region until the 1970s. This population was regarded as a distinct subspecies and assessed as extinct in 2003.
Results of a phylogeographic analysis evinces that the Caspian and Siberian tiger populations shared a common continuous geographic distribution until the early 19th century.
Some Caspian tigers were intermediate in size between Siberian and Bengal tigers.
It was also called Balkhash tiger, Hyrcanian tiger, Turanian tiger, and Mazandaran tiger.
Taxonomy
Felis virgata was a scientific name used by Johann Karl Wilhelm Illiger in 1815 for the greyish tiger in the area surrounding the Caspian Sea.
Tigris septentrionalis was the scientific name proposed by Konstantin Satunin in 1904 for a skull and mounted skins of tigers that were killed in the Lankaran Lowland in the 1860s.
Felis tigris lecoqi and Felis tigris trabata were proposed by Ernst Schwarz in 1916 for tiger skins and skulls from Lop Nur and Ili River areas, respectively.
In 1929, Reginald Innes Pocock subordinated the tiger to the genus Panthera. For several decades, the Caspian tiger was considered a distinct tiger subspecies.
In 1999, the validity of several tiger subspecies was questioned. Most putative subspecies described in the 19th and 20th centuries were distinguished on basis of fur length and colouration, striping patterns and body size, hence characteristics that vary widely within populations. Morphologically, tigers from different regions vary little, and gene flow between populations in those regions is considered to have been possible during the Pleistocene. Therefore, it was proposed to recognize only two tiger subspecies as valid, namely P. t. tigris in mainland Asia, and P. t. sondaica in the Greater Sunda Islands and possibly in Sundaland.
At the start of the 21st century, genetic studies were carried out using 20 tiger bone and tissue samples from museum collections and sequencing at least one segment of five mitochondrial genes. Results revealed a low amount of variability in the mitochondrial DNA in Caspian tigers; and that Caspian and Siberian tigers were remarkably similar, indicating that the Siberian tiger is the genetically closest living relative of the Caspian tiger. Phylogeographic analysis indicates that the common ancestor of Caspian and Siberian tigers colonized Central Asia via the Gansu−Silk Road region from eastern China less than 10,000 years ago, and subsequently traversed eastward to establish the Siberian tiger population in the Russian Far East. The Caspian and Siberian tigers were likely a single contiguous population until the early 19th century, but became isolated from another due to fragmentation and loss of habitat during the Industrial Revolution.
In 2015, morphological, ecological and molecular traits of all putative tiger subspecies were analysed in a combined approach. Results support distinction of the two evolutionary groups continental and Sunda tigers. The authors proposed recognition of only two subspecies, namely P. t. tigris comprising the Bengal, Malayan, Indochinese, South Chinese, Siberian and Caspian tiger populations, and P. t. sondaica comprising the Javan, Bali and Sumatran tiger populations. Tigers in mainland Asia fall into two clades, namely a northern clade formed by the Caspian and Siberian tiger populations, and a southern clade formed by populations in remaining mainland Asia.
In 2017, the Cat Specialist Group revised felid taxonomy and now recognizes the tiger populations in continental Asia as P. t. tigris.
However, a genetic study published in 2018 supported six monophyletic clades, with the Amur and Caspian tigers being distinct from other mainland Asian populations, thus supporting the traditional concept of six living subspecies.
Characteristics
Fur
Photographs of skins of Caspian and Siberian tigers indicate that the main background colour of the Caspian tiger's fur varied and was generally brighter and more uniform than that of the Siberian tiger. The stripes were narrower, fuller and more closely set than those of tigers from Manchuria. The colour of its stripes was a mixture of brown or cinnamon shades. Pure black patterns were invariably found only on head, neck, the middle of the back and at the tip of the tail. Angular patterns at the base of the tail were less developed than those of Far Eastern populations. The contrast between the summer and winter coats was sharp, though not to the same extent as in Far Eastern populations. The winter coat was paler, with less distinct patterns. The summer coat had a similar density and hair length to that of the Bengal tiger, though its stripes were usually narrower, longer and closer set. It had the thickest fur amongst tigers, possibly due its occurrence in the temperate parts of Asia.
Size
Male Caspian tigers had a body length of and weighed ; females measured in head-to-body and weighed . Maximum skull length in males was , while that of females was . Its occiput was broader than of the Bengal tiger.
It ranked among the largest extant cat species, along with the Siberian tiger.
Some individuals attained exceptional sizes. In 1954, a tiger was killed near the Sumbar River in Kopet-Dag, whose stuffed skin was put on display in a museum in Ashgabat. Its head-to-body length was . Its skull had a condylobasal length of about , and zygomatic width of . Its skull length was , hence more than the known maximum of for this population, and slightly exceeding skull length of most Siberian tigers.
In Prishibinske, a tiger was killed in February 1899. Measurements after skinning revealed a body length of between the pegs, plus a long tail, giving it a total length of about . Measurements between the pegs of up to are known. It was said to have been "a tiger of immense proportions" and "no smaller than the local horse breeds." It had rather long fur.
Skull size and shape of Caspian tigers significantly overlap with and are almost indistinguishable from other tiger specimens in mainland Asia.
Distribution and habitat
In the 19th century, tigers occurred in:
the Eastern Anatolia Region, which is considered to have been the westernmost area where tigers occurred. Records are known from the region of Mount Ararat, Şanlıurfa, Şırnak, Siirt and Hakkari Provinces in eastern Turkey; in the Hakkari Province tigers possibly occurred up to the 1990s. The only confirmed record in Iraq dates to 1887 when a tiger was shot near Mosul, which is considered to have been a migrant from southeastern Turkey. There are also claims of historical tiger presence in the area of the Tigris–Euphrates river system in Iraq and Syria.
the extreme southeast of the Caucasus, such as in hilly and lowland forests of the Talysh Mountains, in the Lenkoran Lowlands, in the lowland forests of Göytəpə, Jalilabad, from where tigers moved into the eastern plains of the Trans-Caucasus up to the Don River basin; the Arasbaran and Zangezur Mountains of northwestern Persia.
in the region of the Caspian Sea, where its distribution was patchy and associated with wetlands such as river basins, lake edges and sea shores. In Iran, historical records are known only from along the southern coast of the Caspian Sea and adjacent Alborz Mountains.
Central Asia, such as in southwestern Turkmenia along the Atrek River and its tributaries, and the Sumbar and Chandyr Rivers; in the western and southwestern parts of Kopet-Dag; in the environs of Ashkabad in the northern foothills; in Afghanistan along the upper reaches of Hari-Rud at Herat, and along the jungles in the lower reaches of the river; around Tejen and Murgap and along the Kushka and Kashan rivers; in the Amu Darya basin as far the Aral Sea and along the entire coast of the Aral Sea; along the Syr-Darya to the Fergana Valley as far as Tashkent and the western spur of Talas Alatau; along the Chu and Ili Rivers; all along the southern shore of Lake Balkhash and northwards into the southern Altai Mountains, and to southeastern Transbaikal or Western Siberia in the east. In China, it occurred in the Tarim, Manasi River and Lop Nur basins.
Its former distribution can be approximated by examining the distribution of ungulates in the region. Wild boar was the numerically dominant ungulate in forested habitats, along watercourses, in reed beds and in thickets of the Caspian and Aral Seas. Where watercourses penetrated deep into desert areas, suitable wild pig and tiger habitat was often linear, only a few kilometers wide at most. Red and roe deer occurred in forests around the Black Sea to the western side and around the southern side of the Caspian Sea in a narrow belt of forest cover. Roe deer occurred in forested areas south of Lake Balkhash. Bactrian deer lived in the narrow belt of forest habitat on the southern border of the Aral Sea, and southward along the Syr-Darya and Amu Darya rivers.
Throughout the late Pleistocene and Holocene, the Caspian tiger population was likely connected to the Bengal tiger population through corridors below elevations of in the Hindu Kush, before gene flow was interrupted by humans.
Local extinction
The demise of the Caspian tiger began with the Russian colonisation of Turkestan during the late 19th century.
Its extirpation was caused by several factors:
Tigers were killed by large parties of sportsmen and military personnel who also hunted tiger prey species such as the Bactrian deer, and middle asian Wild boar. This wild pig's range underwent a rapid decline between the middle of the 19th century and the 1930s due to overhunting, natural disasters, and diseases such as swine fever and foot-and-mouth disease, which caused large and rapid die-offs.
The extensive reed beds of tiger habitat were increasingly converted to cropland for planting cotton and other crops that grew well in the rich silt along rivers.
The tiger was already vulnerable due to the restricted nature of its distribution, having been confined to watercourses within the large expanses of desert environment.
Until the early 20th century, the regular Russian army was used to clear predators from forests, around settlements, and potential agricultural lands. Until World War I, about 50 tigers were killed in the forests of Amu Darya and Piandj Rivers each year. High incentives were paid for tiger skins up to 1929. Wild pigs and deer, the prey base of tigers, were decimated by deforestation and subsistence hunting by the increasing human population along the rivers, supported by growing agricultural developments. By 1910, cotton plants were estimated to occupy nearly one-fifth of Turkestan's arable land, with about one half located in the Fergana Valley.
Last sightings
In Iraq, a tiger was killed near Mosul in 1887.
In Georgia, the last known tiger was killed in 1922 near Tbilisi, after taking domestic livestock.
In China, tigers disappeared from the Tarim River basin in Xinjiang in the 1920s.
In Azerbaijan, the last known tiger was killed in 1932; however, tigers were allegedly sighted in later years in the Talysh Mountains.
In Turkey, a pair of tigers was allegedly killed in the area of Selçuk in 1943. Several tiger skins found in the early 1970s near Uludere indicated the presence of a tiger population in eastern Turkey. Questionnaire surveys conducted in this region revealed that one to eight tigers were killed each year until the mid-1980s, and that tigers likely had survived in the region until the early 1990s. Due to lack of interest, in addition to security and safety reasons, no further field surveys were carried out in the area.
In Iran, one of the last known tigers was shot in Golestan National Park in 1953. Another individual was sighted in Golestān Province in 1958.
In Turkmenistan, the last known tiger was killed in January 1954 in the Sumbar River valley in the Kopet-Dag Range.
It reportedly disappeared in the Manasi River basin in the Tian Shan Range west of Ürümqi in the 1960s. The last record from the lower reaches of the Amu Darya river was an unconfirmed observation in 1968 near Nukus in the Aral Sea area. By the early 1970s, tigers disappeared from the river's lower reaches and the Pyzandh Valley in the Turkmen-Uzbek-Afghan border region.
The Piandj River area between Afghanistan and Tajikistan was a stronghold of the Caspian tiger until the late 1960s. The latest sighting of a tiger in the Afghan-Tajik border area dates to 1998 in the Babatag Range. Two tigers were captured in April 1997 in Afghanistan's Laghman Province.
In Kazakhstan, the last Caspian tiger was recorded in 1948, in the environs of the Ili River, the last known stronghold in the region of Lake Balkhash. In May 2006, a Kazakh hunter claimed to have seen a female Caspian tiger with cubs near Lake Balkhash. However, this sighting remains uncertain and unconfirmed.
Behaviour and ecology
No information is available for home ranges of Caspian tigers. In search for prey, they possibly prowled widely and followed migratory ungulates from one pasture to another.
Wild pigs and cervids probably formed their main prey base. In many regions of Central Asia, Bactrian deer and roe deer were important prey species, as well as Caspian red deer and goitered gazelle in Iran; Eurasian golden jackals, jungle cats, locusts, and other small mammals in the lower Amu Darya River area; saigas, wild horses and Persian onagers in the Miankaleh Peninsula; Turkmenian kulans, Mongolian wild asses, and mountain sheep in the Zhana-Darya and around the Aral Sea; and Manchurian wapiti and moose in the area of Lake Baikal. They caught fish in flooded areas and irrigation channels. In winter, they frequently attacked dogs and livestock straying away from herds. They preferred drinking water from rivers, and drank from lakes in seasons when water was less brackish.
Disease
Two tigers in southwestern Tajikistan harbored 5–7 tapeworms (Taenia bubesei) in their small and large intestines.
Conservation
In 1938, Tigrovaya Balka was the first protected area in Tajikistan established in the lower reaches of Vakhsh River between the Panj and Kofarnihon Rivers; it was apparently the last refuge of the Caspian tiger. A tiger was seen there in 1958.
After 1947, tigers were legally protected in the Union of Soviet Socialist Republics.
In Iran, Caspian tigers had been protected since 1957, with heavy fines for shooting. In the early 1970s, biologists from the Department of Environment searched several years for Caspian tigers in the uninhabited areas of Caspian forests, but did not find any evidence of their presence.
In captivity
A tiger from the Caucasus was housed at Berlin Zoo in the late 19th century. A tigress caught in Turkestan was presented to London Zoo on 12 December 1885. DNA from a tiger caught in northern Iran and housed at Moscow Zoo in the 20th century was used in the genetic test that established the Caspian tiger's close genetic relationship with the Siberian tiger. This tigress lived from 1924 to 1942 and was presented to the Soviet ambassador in Iran. Another tigress kept at Tierpark Hagenbeck in Hamburg between 1955 and 1960 was probably the last Caspian tiger in captivity. An individual was born in Brookfield Zoo Chicago on 7 May 1935 and was still living on 1 January 1948.
Reintroduction project
Stimulated by recent findings that the Siberian tiger is the closest relative of the Caspian tiger, discussions started as to whether the Siberian tiger could be appropriate for reintroduction into a safe place in Central Asia, where the Caspian tiger once roamed. The Amu Darya delta was suggested as a potential site for such a project. A feasibility study was initiated to investigate if the area is suitable, and if such an initiative would receive support from relevant decision makers. A viable tiger population of about 100 animals would require at least of large tracts of contiguous habitat, with rich prey populations. Such habitat is not currently available, and cannot be provided in the short term. The proposed region is therefore unsuitable for the reintroduction, at least at the current stage.
While the restoration of the Caspian tiger has stimulated discussions, the locations for the tiger have yet to become fully involved in the planning. But through preliminary ecological surveys it has been revealed that some small populated areas of Central Asia have preserved natural habitat suitable for tigers.
In culture
In the Roman Empire, tigers and other large animals imported from Africa and Asia were used during gladiatorial games.
In the Taurus Mountains, stone traps were used to capture leopards and tigers.
In the Fables of Pilpay, the tiger is described as furious and avid to rule over wilderness.
The babr (, tiger) features in Persian and Central Asian culture. The name "Babr Mazandaran" is sometimes given to a prominent wrestler. A Syrian mosaic in Palmyra depicts the Sassanids as tigers, possibly commemorating the victory of the Palmyrene King Odaenathus over Shapur I. The inscription on the mosaic conceals an earlier one that read: (Mrn), which is a title used by Odaenathus. It possibly celebrates Odaenathus' victory over the Persians, the archer representing Odaenathus and the tigers the Persians; Odaenathus is about to be crowned with victory by the eagle flying above him.
| Biology and health sciences | Felines | Animals |
1458081 | https://en.wikipedia.org/wiki/Yield%20%28chemistry%29 | Yield (chemistry) | In chemistry, yield, also known as reaction yield or chemical yield, refers to the amount of product obtained in a chemical reaction. Yield is one of the primary factors that scientists must consider in organic and inorganic chemical synthesis processes. In chemical reaction engineering, "yield", "conversion" and "selectivity" are terms used to describe ratios of how much of a reactant was consumed (conversion), how much desired product was formed (yield) in relation to the undesired product (selectivity), represented as X, Y, and S.
The term yield also plays an important role in analytical chemistry, as individual compounds are recovered in purification processes in a range from quantitative yield (100 %) to low yield (< 50 %).
Definitions
In chemical reaction engineering, "yield", "conversion" and "selectivity" are terms used to describe ratios of how much of a reactant has reacted—conversion, how much of a desired product was formed—yield, and how much desired product was formed in ratio to the undesired product—selectivity, represented as X, S, and Y.
According to the Elements of Chemical Reaction Engineering manual, yield refers to the amount of a specific product formed per mole of reactant consumed. In chemistry, mole is used to describe quantities of reactants and products in chemical reactions.
The Compendium of Chemical Terminology defined yield as the "ratio expressing the efficiency of a mass conversion process. The yield coefficient is defined as the amount of cell mass (kg) or product formed (kg,mol) related to the consumed substrate (carbon or nitrogen source or oxygen in kg or moles) or to the intracellular ATP production (moles)."
In the section "Calculations of yields in the monitoring of reactions" in the 1996 4th edition of Vogel's Textbook of Practical Organic Chemistry (1978), the authors write that, "theoretical yield in an organic reaction is the weight of product which would be obtained if the reaction has proceeded to completion according to the chemical equation. The yield is the weight of the pure product which is isolated from the reaction." In 'the 1996 edition of Vogel's Textbook, percentage yield is expressed as,
According to the 1996 edition of Vogel's Textbook, yields close to 100% are called quantitative, yields above 90% are called excellent, yields above 80% are very good, yields above 70% are good, yields above 50% are fair, and yields below 40% are called poor. In their 2002 publication, Petrucci, Harwood, and Herring wrote that Vogel's Textbook names were arbitrary, and not universally accepted, and depending on the nature of the reaction in question, these expectations may be unrealistically high. Yields may appear to be 100% or above when products are impure, as the measured weight of the product will include the weight of any impurities.
In their 2016 laboratory manual, Experimental Organic Chemistry, the authors described the "reaction yield" or "absolute yield" of a chemical reaction as the "amount of pure and dry product yielded in a reaction". They wrote that knowing the stoichiometry of a chemical reaction—the numbers and types of atoms in the reactants and products, in a balanced equation "make it possible to compare different elements through stoichiometric factors." Ratios obtained by these quantitative relationships are useful in data analysis.
Theoretical, actual, and percent yields
The percent yield is a comparison between the actual yield—which is the weight of the intended product of a chemical reaction in a laboratory setting—and the theoretical yield—the measurement of pure intended isolated product, based on the chemical equation of a flawless chemical reaction, and is defined as,
The ideal relationship between products and reactants in a chemical reaction can be obtained by using a chemical reaction equation. Stoichiometry is used to run calculations about chemical reactions, for example, the stoichiometric mole ratio between reactants and products. The stoichiometry of a chemical reaction is based on chemical formulas and equations that provide the quantitative relation between the number of moles of various products and reactants, including yields. Stoichiometric equations are used to determine the limiting reagent or reactant—the reactant that is completely consumed in a reaction. The limiting reagent determines the theoretical yield—the relative quantity of moles of reactants and the product formed in a chemical reaction. Other reactants are said to be present in excess. The actual yield—the quantity physically obtained from a chemical reaction conducted in a laboratory—is often less than the theoretical yield. The theoretical yield is what would be obtained if all of the limiting reagent reacted to give the product in question. A more accurate yield is measured based on how much product was actually produced versus how much could be produced. The ratio of the theoretical yield and the actual yield results in a percent yield.
When more than one reactant participates in a reaction, the yield is usually calculated based on the amount of the limiting reactant, whose amount is less than stoichiometrically equivalent (or just equivalent) to the amounts of all other reactants present. Other reagents present in amounts greater than required to react with all the limiting reagent present are considered excess. As a result, the yield should not be automatically taken as a measure for reaction efficiency.
In their 1992 publication General Chemistry, Whitten, Gailey, and Davis described the theoretical yield as the amount predicted by a stoichiometric calculation based on the number of moles of all reactants present. This calculation assumes that only one reaction occurs and that the limiting reactant reacts completely.
According to Whitten, the actual yield is always smaller (the percent yield is less than 100%), often very much so, for several reasons. As a result, many reactions are incomplete and the reactants are not completely converted to products. If a reverse reaction occurs, the final state contains both reactants and products in a state of chemical equilibrium. Two or more reactions may occur simultaneously, so that some reactant is converted to undesired side products. Losses occur in the separation and purification of the desired product from the reaction mixture. Impurities are present in the starting material which do not react to give desired product.
Example
This is an example of an esterification reaction where one molecule acetic acid (also called ethanoic acid) reacts with one molecule ethanol, yielding one molecule ethyl acetate (a bimolecular second-order reaction of the type A + B → C):
120 g acetic acid (60 g/mol, 2.0 mol) was reacted with 230 g ethanol (46 g/mol, 5.0 mol), yielding 132 g ethyl acetate (88 g/mol, 1.5 mol). The yield was 75%.
The molar amount of the reactants is calculated from the weights (acetic acid: 120 g ÷ 60 g/mol = 2.0 mol; ethanol: 230 g ÷ 46 g/mol = 5.0 mol).
Ethanol is used in a 2.5-fold excess (5.0 mol ÷ 2.0 mol).
The theoretical molar yield is 2.0 mol (the molar amount of the limiting compound, acetic acid).
The molar yield of the product is calculated from its weight (132 g ÷ 88 g/mol = 1.5 mol).
The % yield is calculated from the actual molar yield and the theoretical molar yield (1.5 mol ÷ 2.0 mol × 100% = 75%).
Purification of products
In his 2016 Handbook of Synthetic Organic Chemistry, Michael Pirrung wrote that yield is one of the primary factors synthetic chemists must consider in evaluating a synthetic method or a particular transformation in "multistep syntheses." He wrote that a yield based on recovered starting material (BRSM) or (BORSM) does not provide the theoretical yield or the "100% of the amount of product calculated", that is necessary in order to take the next step in the multistep systhesis.
Purification steps always lower the yield, through losses incurred during the transfer of material between reaction vessels and purification apparatus or imperfect separation of the product from impurities, which may necessitate the discarding of fractions deemed insufficiently pure. The yield of the product measured after purification (typically to >95% spectroscopic purity, or to sufficient purity to pass combustion analysis) is called the isolated yield of the reaction.
Internal standard yield
Yields can also be calculated by measuring the amount of product formed (typically in the crude, unpurified reaction mixture) relative to a known amount of an added internal standard, using techniques like Gas chromatography (GC), High-performance liquid chromatography, or Nuclear magnetic resonance spectroscopy (NMR spectroscopy) or magnetic resonance spectroscopy (MRS). A yield determined using this approach is known as an internal standard yield. Yields are typically obtained in this manner to accurately determine the quantity of product produced by a reaction, irrespective of potential isolation problems. Additionally, they can be useful when isolation of the product is challenging or tedious, or when the rapid determination of an approximate yield is desired. Unless otherwise indicated, yields reported in the synthetic organic and inorganic chemistry literature refer to isolated yields, which better reflect the amount of pure product one is likely to obtain under the reported conditions, upon repeating the experimental procedure.
Reporting of yields
In their 2010 Synlett article, Martina Wernerova and organic chemist, Tomáš Hudlický, raised concerns about inaccurate reporting of yields, and offered solutions—including the proper characterization of compounds. After performing careful control experiments, Wernerova and Hudlický said that each physical manipulation (including extraction/washing, drying over desiccant, filtration, and column chromatography) results in a loss of yield of about 2%. Thus, isolated yields measured after standard aqueous workup and chromatographic purification should seldom exceed 94%. They called this phenomenon "yield inflation" and said that yield inflation had gradually crept upward in recent decades in chemistry literature. They attributed yield inflation to careless measurement of yield on reactions conducted on small scale, wishful thinking and a desire to report higher numbers for publication purposes.
| Physical sciences | Reaction | Chemistry |
1458550 | https://en.wikipedia.org/wiki/Depersonalization-derealization%20disorder | Depersonalization-derealization disorder | Depersonalization-derealization disorder (DPDR, DDD) is a mental disorder in which the person has persistent or recurrent feelings of depersonalization and/or derealization. Depersonalization is described as feeling disconnected or detached from one's self. Individuals may report feeling as if they are an outside observer of their own thoughts or body, and often report feeling a loss of control over their thoughts or actions. Derealization is described as detachment from one's surroundings. Individuals experiencing derealization may report perceiving the world around them as foggy, dreamlike, surreal, and/or visually distorted.
Depersonalization-derealization disorder is thought to be caused largely by interpersonal trauma such as early childhood abuse. Adverse childhood experiences, specifically emotional abuse and neglect have been linked to the development of depersonalization symptoms. Feelings of depersonalization and derealization are common from significant stress or panic attacks. Individuals may remain in a depersonalized state for the duration of a typical panic attack. However, in some cases, the dissociated state may last for hours, days, weeks, or even months at a time. In rare cases, symptoms of a single episode can last for years.
Diagnostic criteria for depersonalization-derealization disorder includes persistent or recurrent feelings of detachment from one's mental or bodily processes or from one's surroundings. A diagnosis is made when the dissociation is persistent, interferes with the social or occupational functions of daily life, and/or causes marked distress in the patient.
While depersonalization-derealization disorder was once considered rare, lifetime experiences with it occur in about 1–2% of the general population. The chronic form of the disorder has a reported prevalence of 0.8 to 1.9%. While brief episodes of depersonalization or derealization can be common in the general population, the disorder is only diagnosed when these symptoms cause substantial distress or impair social, occupational, or other important areas of functioning.
Signs and symptoms
The core symptoms of depersonalization-derealization disorder are the subjective experience of "unreality in one's self", or detachment from one's surroundings. People who are diagnosed with depersonalization also often experience an urge to question and think critically about the nature of reality and existence.
Individuals with depersonalization describe feeling disconnected from their physicality; feeling as if they are not completely occupying their own body; feeling as if their speech or physical movements are out of their control; feeling detached from their own thoughts or emotions; and experiencing themselves and their lives from a distance. While depersonalization involves detachment from one's self, individuals with derealization feel detached from their surroundings, as if the world around them is foggy, dreamlike, or visually distorted. Individuals with the disorder commonly describe a feeling as though time is passing them by and they are not in the notion of the present. In some cases, individuals may be unable to accept their reflection as their own, or they may have out-of-body experiences. Additionally some individuals experience difficulty concentrating and problems with memory retrieval. These individuals sometimes lack the "feeling" of a memory where they are able to recall a memory but feel as if they did not personally experience it. These experiences which strike at the core of a person's identity and consciousness may cause a person to feel uneasy or anxious. The inner turmoil created by the disorder can also result in depression.
First experiences with depersonalization may be frightening, with patients fearing loss of control, dissociation from the rest of society and functional impairment. The majority of people with depersonalization-derealization disorder misinterpret the symptoms, thinking that they are signs of serious psychosis or brain dysfunction. This commonly leads to an increase of anxiety and obsession, which contributes to the worsening of symptoms.
Factors that tend to diminish symptoms are comforting personal interactions, intense physical or emotional stimulation, and relaxation. Distracting oneself (by engaging in conversation or watching a movie, for example) may also provide temporary relief. Some other factors that are identified as relieving symptom severity are diet or exercise, while alcohol and fatigue are listed by some as worsening their symptoms.
Occasional, brief moments of mild depersonalization can be experienced by many members of the general population; however, depersonalization-derealization disorder occurs when these feelings are strong, severe, persistent, or recurrent and when these feelings interfere with daily functioning. DPDR is most commonly experienced as chronic and continuous. However, for a minority who have DPDR as an episodic condition, duration of these episodes is highly variable with some lasting as long as several weeks.
Causes
The exact cause of depersonalization is unknown, although biopsychosocial correlations and triggers have been identified. It has been thought that depersonalization can be caused by a biological response to dangerous or life-threatening situations which causes heightened senses and emotional numbing.
Psychosocial
There is growing evidence linking physical and sexual abuse in childhood with the development of dissociative disorders. Childhood interpersonal trauma – emotional abuse in particular – is a significant predictor of a diagnosis of DPDR. Compared to other types of childhood trauma, emotional abuse has been found to be the most significant predictor both of a diagnosis of depersonalization disorder and of depersonalization scores, but not of general dissociation scores. Some studies suggest that greater emotional abuse and lower physical abuse predict depersonalization in adult women with post-traumatic stress disorder (PTSD). Patients with high interpersonal abuse histories (HIA) show significantly higher scores on the Cambridge Depersonalization Scale, when compared to a control group. Earlier age of abuse, increased duration and parental abuse tend to correlate with severity of dissociative symptoms. Besides traumatic experiences, other common precipitators of the disorder include severe stress, major depressive disorder or panic attacks. People who live in highly individualistic cultures may be more vulnerable to depersonalization due to a hypersensitivity towards threats and fears of losing control.
A 2010 study found evidence that some users participating in virtual reality (VR) may be more likely to experience dissociation after use. Users reportedly experienced higher levels of a lessened sense of presence in reality after exposure to VR. However, it was noted that the effects of exposure were likely to rapidly disappear after returning to objective reality. Additionally, individuals who reported higher preexisting dissociation levels as well being more easily immersed or absorbed in imagination overall were found to be linked to higher increases in dissociative symptoms after the VR exposure. This study offered evidence towards a link between imaginative processes of the brain and dissociative experiences.
Neurobiology
There is converging evidence that the prefrontal cortex may inhibit neural circuits that normally form the basis of emotional experience. In an fMRI study of DPDR patients, emotionally aversive scenes activated the right ventral prefrontal cortex. Participants demonstrated a reduced neural response in emotion-sensitive regions, as well as an increased response in regions associated with emotional regulation. In a similar test of emotional memory, depersonalization disorder patients did not process emotionally salient material in the same way as did healthy controls. In a test of skin conductance responses to unpleasant stimuli, the subjects showed a selective inhibitory mechanism on emotional processing.
Studies are beginning to show that the temporoparietal junction has a role in multisensory integration, embodiment, and self-other distinction. Several studies analyzing brain MRI findings from DPDR patients found decreased cortical thickness in the right middle temporal gyrus, reduction in grey matter volume in the right caudate, thalamus, and occipital gyri, as well as lower white matter integrity in the left temporal and right temporoparietal regions. However, no structural changes in the amygdala were observed.
A PET scan found functional abnormalities in the visual, auditory, and somatosensory cortex, as well as in areas responsible for an integrated body schema.
One study examining EEG readings found frontal alpha wave overactivation and increased theta activity waves in the temporal region of the left hemisphere.
It is unclear whether genetics plays a role; however, there are many neurochemical and hormonal changes in individuals with depersonalization disorder. DPDR may be associated with dysregulation of the hypothalamic-pituitary-adrenal axis, the area of the brain involved in the "fight-or-flight" response. Patients demonstrate abnormal cortisol levels and basal activity. Studies found that patients with DPDR could be distinguished from patients with clinical depression and posttraumatic stress disorder.
The vestibular system may also play a role in DPDR. The vestibular system helps control balance, spatial orientation, motor coordination, but also plays a role in self-awareness. Disruption to this system can potentially cause a feeling of detachment from surroundings. Several studies have shown that patients with peripheral vestibular disease are also more likely to have dissociative symptoms when compared to healthy individuals.
Dissociative symptoms are sometimes described by those with neurological diseases, such as amyotrophic lateral sclerosis, Alzheimer's, multiple sclerosis (MS), etc., that directly affect brain tissue.
Diagnosis
Assessment
Diagnosis is based on the self-reported experiences of the person followed by a clinical assessment. Psychiatric assessment includes a psychiatric history and some form of mental status examination. Since some medical and psychiatric conditions mimic the symptoms of DPDR, clinicians must differentiate between and rule out the following to establish a precise diagnosis: temporal lobe epilepsy, panic disorder, acute stress disorder, schizophrenia, migraine, drug use, brain tumor or lesion. No laboratory test for depersonalization-derealization disorder currently exists. As patients with dissociative disorders likely experienced intense trauma in the past, concomitant dissociative disorders should be considered in patients diagnosed with a stress disorder (i.e. PTSD or acute stress disorder).
The diagnosis of depersonalization disorder can be made with the use of the following interviews and scales:
The Structured Clinical Interview for DSM-IV Dissociative Disorders (SCID-D) is widely used, especially in research settings. This interview takes about 30 minutes to 1.5 hours, depending on individual's experiences.
The Dissociative Experiences Scale (DES) is a simple, quick, self-administered questionnaire that has been widely used to measure dissociative symptoms. It has been used in hundreds of dissociative studies, and can detect depersonalization and derealization experiences.
The Dissociative Disorders Interview Schedule (DDIS) is a highly structured interview which makes DSM-IV diagnoses of somatization disorder, borderline personality disorder and major depressive disorder, as well as all the dissociative disorders. It inquires about positive symptoms of schizophrenia, secondary features of dissociative identity disorder, extrasensory experiences, substance abuse and other items relevant to the dissociative disorders. The DDIS can usually be administered in 30–45 minutes.
The Cambridge Depersonalization Scale (CDS) is a method for determining the severity of depersonalization disorder. It has been proven and accepted as a valid tool for the diagnosis of depersonalization disorder in a clinical setting. It is also used in a clinical setting to differentiate minor episodes of depersonalization from actual symptoms of the disorder. Due to the success of the CDS, a group of Japanese researchers underwent the effort to translate the CDS into the J-CDS or the Japanese Cambridge Depersonalization Scale. Through clinical trials, the Japanese research team successfully tested their scale and determined its accuracy. One limitation is that the scale does not allow for the differentiation between past and present episodes of depersonalization. It may be difficult for the individual to describe the duration of a depersonalization episode, and thus the scale may lack accuracy. The project was conducted in the hope that it would stimulate further scientific investigations into depersonalization disorder.
Diagnostic and Statistical Manual of Mental Disorders, 5th Edition (DSM-5)
In the DSM-5, the word "derealization" was added to "depersonalization disorder" and renamed "depersonalization/derealization disorder" ("DPDR"). It remains classified as a dissociative disorder.
Patients must meet the following criteria to be diagnosed per the DSM-5:
Presence of persistent/recurrent episodes of depersonalization/derealization
Ability to distinguish between reality and dissociation during an episode (i.e. patient is aware of a perceptual disturbance)
Symptoms are severe enough to interfere with social, occupational, or other areas of functioning
Symptoms are not due to a substance or medication
Symptoms are not due to another psychiatric disorder
International Classification of Diseases 11th Revision (ICD-11)
The ICD-11 has relisted DPDR as a disorder rather than a syndrome as previously, and has also reclassified it as a dissociative disorder from its previous listing as a neurotic disorder. The description used in the ICD-11 is similar to the criteria found in the DSM-5. Individuals with DPDR are described as having persistent/recurrent symptoms of depersonalization/derealization, have intact reality testing, and symptoms are not better explained by another psychiatric/neural disorder, substance, medication, or head trauma. Symptoms are severe enough to cause distress or impairment in functioning.
Differential diagnoses
DPDR differentials include neurologic and psychiatric conditions as well as side effects from psychoactive substances or medications.
Neurologic
Seizures
Brain tumor
Post-concussion syndrome
Metabolic abnormalities
Migraines
Vertigo
Meniere's disease
Visual snow syndrome
Psychiatric
Panic attack
Phobias
Post-traumatic stress disorder
Acute stress disorder
Depression
Bipolar disorder
Schizophrenia
Borderline personality disorder
Other dissociative disorders
Dissociative identity disorder
Consequence of psychoactive substance use
Marijuana
Hallucinogens
MDMA
Ketamine
Hallucinogen persisting perception disorder
Prevention
Depersonalization-derealization disorder may be prevented by connecting children who have been abused with professional mental health help. Some trauma specialists strongly advocate for increasing inquiry into information about children's trauma history and exposure to violence, since the majority of people (about 80%) responsible for child maltreatment are the child's own caregivers. Trauma-specific intervention for children may be useful in preventing future symptoms.
Treatment
Treatment of DPDR is often difficult and refractory. Some clinicians speculate that this could be due to a delay in diagnosis by which point symptoms tend to be constant and less responsive to treatment. Additionally, symptoms tend to overlap with other diagnoses. Some results have been promising, but are hard to evaluate with confidence due to the small size of trials. However, recognizing and diagnosing the condition may in itself have therapeutic benefits, considering many patients express their problems as baffling and unique to them, but are not, in fact, and are recognized and described by psychiatry. However, symptoms are often transient and can remit on their own without treatment.
Treatment is primarily pharmacological. Self-hypnosis training can be helpful and entails training patients to induce dissociative symptoms and respond in an alternative manner. Psychoeducation involves counseling regarding the disorder, reassurance, and emphasis on DPDR as a perceptual disturbance rather than a true physical experience. Clinical pharmacotherapy research continues to explore a number of possible options, including selective serotonin reuptake inhibitors (SSRI), benzodiazepines, stimulants and opioid antagonists (ex: naltrexone).
Cognitive behavioral therapy
An open study of cognitive behavioral therapy has aimed to help patients reinterpret their symptoms in a nonthreatening way, leading to an improvement on several standardized measures. A standardized treatment for DPDR based on cognitive behavioral principles was published in the Netherlands in 2011.
Medications
Tentative evidence supports the use of opioid antagonists (naloxone) and other medications like benzodiazepines or methylphenidate. Evidence suggests the beneficial use of lamotrigine adjunct to an SSRI but not as monotherapy.
A combination of an SSRI and a benzodiazepine has been proposed to be useful for DPDR patients with anxiety.
Modafinil used alone has been reported to be effective in a subgroup of individuals with depersonalization disorder (those who have attentional impairments, under-arousal and hypersomnia). However, clinical trials have not been conducted.
Repetitive transcranial magnetic stimulation (rTMS)
Some studies have found repetitive transcranial magnetic stimulation (rTMS) to be helpful. One study examined 12 patients with DPDR that were treated with right temporoparietal junction (TPJ) rTMS and found that 50% showed improvement after three weeks of treatment. Five of the participants received an additional three weeks of treatment and reported overall a 68% improvement in their symptoms. Treating patients with rTMS specifically at the TPJ may be an alternative treatment.
Prognosis
Michal et al. (2016) analyzed a case series on 223 patients suffering from DPDR and agreed that the condition tended to be long-lasting. However, while no medication has been confirmed to successfully treat the condition, psychotherapy might help. In some cases, recovery can take place organically, without formal treatment.
Epidemiology
Men and women are diagnosed in equal numbers with depersonalization disorder. A 1991 study on a sample from Winnipeg, Manitoba estimated the prevalence of depersonalization disorder at 2.4% of the population. A 2008 review of several studies estimated the prevalence between 0.8% and 1.9%. This disorder is episodic in only one-third of individuals, with each episode lasting from hours to months at a time. Depersonalization can begin episodically, and later become continuous at constant or varying intensity.
Onset is typically during adolescence, although some patients report being depersonalized as long as they can remember, and a small minority report a later onset (by age 40). According to the DSM-5-TR, less than 20% of patients with the disorder first experience symptoms after age 20 years; 80% or more have their onset in the first 2 decades of life - childhood and adolescence. The onset can be acute or insidious in nature. With acute onset, some individuals remember the exact time and place of their first experience of depersonalization and/or derealization. This may follow a prolonged period of severe stress, a traumatic event, or an episode of another mental illness. Insidious onset may reach back as far as can be remembered (early childhood), or it may begin with smaller episodes of lesser severity that become gradually more intense and more disabling. Some patients report persistent depersonalization and/or derealization throughout the day, nearly everyday.
Relation to other psychiatric disorders
Depersonalization exists as both a primary and secondary phenomenon. The most common comorbid disorders are depression and anxiety, although cases of depersonalization disorder without symptoms of either do exist. Comorbid obsessive/compulsive behaviors may exist as attempts to deal with depersonalization, such as checking whether symptoms have changed and avoiding behavioral and cognitive factors that exacerbate symptoms. Many people with personality disorders such as schizoid personality disorder, schizotypal personality disorder, and borderline personality disorder will have experiences of depersonalization. Patients with complex dissociative disorders, including dissociative identity disorder, experience high levels of depersonalization and derealization.
History
The word depersonalization itself was first used by Henri Frédéric Amiel in The Journal Intime. The 8 July 1880 entry reads:
Depersonalization was first used as a clinical term by Ludovic Dugas in 1898 to refer to "a state in which there is the feeling or sensation that thoughts and acts elude the self and become strange; there is an alienation of personality – in other words a depersonalization". This description refers to personalization as a psychical synthesis of attribution of states to the self.
Early theories of the cause of depersonalization focused on sensory impairment. Maurice Krishaber proposed depersonalization was the result of pathological changes to the body's sensory modalities which lead to experiences of "self-strangeness" and the description of one patient who "feels that he is no longer himself". One of Carl Wernicke's students suggested all sensations were composed of a sensory component and a related muscular sensation that came from the movement itself and served to guide the sensory apparatus to the stimulus. In depersonalized patients, these two components were not synchronized, and the myogenic sensation failed to reach consciousness. The sensory hypothesis was challenged by others who suggested that patient complaints were being taken too literally and that some descriptions were metaphors – attempts to describe experiences that are difficult to articulate in words. Pierre Janet approached the theory by pointing out his patients with clear sensory pathology did not complain of symptoms of unreality, and that those who have depersonalization were normal from a sensory viewpoint.
Psychodynamic theory formed the basis for the conceptualization of dissociation as a defense mechanism. Within this framework, depersonalization is understood as a defense against a variety of negative feelings, conflicts, or experiences. Sigmund Freud himself experienced fleeting derealization when visiting the Acropolis in person; having read about it for years and knowing it existed, seeing the real thing was overwhelming and proved difficult for him to perceive it as real. Freudian theory is the basis for the description of depersonalization as a dissociative reaction, placed within the category of psychoneurotic disorders, in the first two editions of the Diagnostic and Statistical Manual of Mental Disorders.
It can be argued that because depersonalization and derealization are both impairments to one's ability to perceive reality, they are merely two facets of the same disorder. Depersonalization also differs from delusion in the sense that the patient is able to differentiate between reality and the symptoms they may experience. The ability to sense that something is unreal is maintained when experiencing symptoms of the disorder. The problem with properly defining depersonalization also lies within the understanding of what reality actually is. In order to comprehend the nature of reality we must incorporate all the subjective experiences throughout and thus the problem of obtaining an objective definition is brought about again.
Society and culture
Depersonalization disorder has appeared in a variety of media. The director of the autobiographical documentary Tarnation, Jonathan Caouette, had depersonalization disorder. The screenwriter for the 2007 film Numb had depersonalization disorder, as does the film's protagonist played by Matthew Perry. Norwegian painter Edvard Munch's famous masterpiece The Scream may have been inspired by depersonalization disorder. In Glen Hirshberg's novel The Snowman's Children, main female plot characters throughout the book had a condition that is revealed to be depersonalization disorder. Suzanne Segal had an episode in her 20s that was diagnosed by several psychologists as depersonalization disorder, though Segal herself interpreted it through the lens of Buddhism as a spiritual experience, commonly known as "Satori" or "Samadhi". The song "Is Happiness Just a Word?" by hip hop artist Vinnie Paz describes his struggle with depersonalization disorder. Adam Duritz, of the band Counting Crows, has often spoken about his diagnosis of depersonalization disorder.
| Biology and health sciences | Mental disorders | Health |
1460126 | https://en.wikipedia.org/wiki/Chromatic%20polynomial | Chromatic polynomial | The chromatic polynomial is a graph polynomial studied in algebraic graph theory, a branch of mathematics. It counts the number of graph colorings as a function of the number of colors and was originally defined by George David Birkhoff to study the four color problem. It was generalised to the Tutte polynomial by Hassler Whitney and W. T. Tutte, linking it to the Potts model of statistical physics.
History
George David Birkhoff introduced the chromatic polynomial in 1912, defining it only for planar graphs, in an attempt to prove the four color theorem. If denotes the number of proper colorings of G with k colors then one could establish the four color theorem by showing for all planar graphs G. In this way he hoped to apply the powerful tools of analysis and algebra for studying the roots of polynomials to the combinatorial coloring problem.
Hassler Whitney generalised Birkhoff’s polynomial from the planar case to general graphs in 1932. In 1968, Ronald C. Read asked which polynomials are the chromatic polynomials of some graph, a question that remains open, and introduced the concept of chromatically equivalent graphs. Today, chromatic polynomials are one of the central objects of algebraic graph theory.
Definition
For a graph G, counts the number of its (proper) vertex k-colorings.
Other commonly used notations include , , or .
There is a unique polynomial which evaluated at any integer k ≥ 0 coincides with ; it is called the chromatic polynomial of G.
For example, to color the path graph on 3 vertices with k colors, one may choose any of the k colors for the first vertex, any of the remaining colors for the second vertex, and lastly for the third vertex, any of the colors that are different from the second vertex's choice.
Therefore, is the number of k-colorings of .
For a variable x (not necessarily integer), we thus have .
(Colorings which differ only by permuting colors or by automorphisms of G are still counted as different.)
Deletion–contraction
The fact that the number of k-colorings is a polynomial in k follows from a recurrence relation called the deletion–contraction recurrence or Fundamental Reduction Theorem. It is based on edge contraction: for a pair of vertices and the graph is obtained by merging the two vertices and removing any edges between them.
If and are adjacent in G, let denote the graph obtained by removing the edge .
Then the numbers of k-colorings of these graphs satisfy:
Equivalently, if and are not adjacent in G and is the graph with the edge added, then
This follows from the observation that every k-coloring of G either gives different colors to and , or the same colors. In the first case this gives a (proper) k-coloring of , while in the second case it gives a coloring of .
Conversely, every k-coloring of G can be uniquely obtained from a k-coloring of or (if and are not adjacent in G).
The chromatic polynomial can hence be recursively defined as
for the edgeless graph on n vertices, and
for a graph G with an edge (arbitrarily chosen).
Since the number of k-colorings of the edgeless graph is indeed , it follows by induction on the number of edges that for all G, the polynomial coincides with the number of k-colorings at every integer point x = k. In particular, the chromatic polynomial is the unique interpolating polynomial of degree at most n through the points
Tutte’s curiosity about which other graph invariants satisfied such recurrences led him to discover a bivariate generalization of the chromatic polynomial, the Tutte polynomial .
Examples
Properties
For fixed G on n vertices, the chromatic polynomial is a monic polynomial of degree exactly n, with integer coefficients.
The chromatic polynomial includes at least as much information about the colorability of G as does the chromatic number. Indeed, the chromatic number is the smallest positive integer that is not a zero of the chromatic polynomial,
The polynomial evaluated at , that is , yields times the number of acyclic orientations of G.
The derivative evaluated at 1, equals the chromatic invariant up to sign.
If G has n vertices and c components , then
The coefficients of are zeros.
The coefficients of are all non-zero and alternate in signs.
The coefficient of is 1 (the polynomial is monic).
The coefficient of is
We prove this via induction on the number of edges on a simple graph G with vertices and edges. When , G is an empty graph. Hence per definition . So the coefficient of is , which implies the statement is true for an empty graph. When , as in G has just a single edge, . Thus coefficient of is . So the statement holds for k = 1. Using strong induction assume the statement is true for . Let G have edges. By the contraction-deletion principle, Let and Hence .Since is obtained from G by removal of just one edge e, , so and thus the statement is true for k.
The coefficient of is times the number of acyclic orientations that have a unique sink, at a specified, arbitrarily chosen vertex.
The absolute values of coefficients of every chromatic polynomial form a log-concave sequence.
The last property is generalized by the fact that if G is a k-clique-sum of and (i.e., a graph obtained by gluing the two at a clique on k vertices), then
A graph G with n vertices is a tree if and only if
Chromatic equivalence
Two graphs are said to be chromatically equivalent if they have the same chromatic polynomial. Isomorphic graphs have the same chromatic polynomial, but non-isomorphic graphs can be chromatically equivalent. For example, all trees on n vertices have the same chromatic polynomial.
In particular, is the chromatic polynomial of both the claw graph and the path graph on 4 vertices.
A graph is chromatically unique if it is determined by its chromatic polynomial, up to isomorphism. In other words, G is chromatically unique, then would imply that G and H are isomorphic.
All cycle graphs are chromatically unique.
Chromatic roots
A root (or zero) of a chromatic polynomial, called a “chromatic root”, is a value x where . Chromatic roots have been very well studied, in fact, Birkhoff’s original motivation for defining the chromatic polynomial was to show that for planar graphs, for x ≥ 4. This would have established the four color theorem.
No graph can be 0-colored, so 0 is always a chromatic root. Only edgeless graphs can be 1-colored, so 1 is a chromatic root of every graph with at least one edge. On the other hand, except for these two points, no graph can have a chromatic root at a real number smaller than or equal to 32/27. A result of Tutte connects the golden ratio with the study of chromatic roots, showing that chromatic roots exist very close to :
If is a planar triangulation of a sphere then
While the real line thus has large parts that contain no chromatic roots for any graph, every point in the complex plane is arbitrarily close to a chromatic root in the sense that there exists an infinite family of graphs whose chromatic roots are dense in the complex plane.
Colorings using all colors
For a graph G on n vertices, let denote the number of colorings using exactly k colors up to renaming colors (so colorings that can be obtained from one another by permuting colors are counted as one; colorings obtained by automorphisms of G are still counted separately).
In other words, counts the number of partitions of the vertex set into k (non-empty) independent sets.
Then counts the number of colorings using exactly k colors (with distinguishable colors).
For an integer x, all x-colorings of G can be uniquely obtained by choosing an integer k ≤ x, choosing k colors to be used out of x available, and a coloring using exactly those k (distinguishable) colors.
Therefore:
where denotes the falling factorial.
Thus the numbers are the coefficients of the polynomial in the basis of falling factorials.
Let be the k-th coefficient of in the standard basis , that is:
Stirling numbers give a change of basis between the standard basis and the basis of falling factorials.
This implies:
and
Categorification
The chromatic polynomial is categorified by a homology theory closely related to Khovanov homology.
Algorithms
Computational problems associated with the chromatic polynomial include
finding the chromatic polynomial of a given graph G;
evaluating at a fixed x for given G.
The first problem is more general because if we knew the coefficients of we could evaluate it at any point in polynomial time because the degree is n. The difficulty of the second type of problem depends strongly on the value of x and has been intensively studied in computational complexity. When x is a natural number, this problem is normally viewed as computing the number of x-colorings of a given graph. For example, this includes the problem #3-coloring of counting the number of 3-colorings, a canonical problem in the study of complexity of counting, complete for the counting class #P.
Efficient algorithms
For some basic graph classes, closed formulas for the chromatic polynomial are known. For instance this is true for trees and cliques, as listed in the table above.
Polynomial time algorithms are known for computing the chromatic polynomial for wider classes of graphs, including chordal graphs and graphs of bounded clique-width. The latter class includes cographs and graphs of bounded tree-width, such as outerplanar graphs.
Deletion–contraction
The deletion-contraction recurrence gives a way of computing the chromatic polynomial, called the deletion–contraction algorithm. In the first form (with a minus), the recurrence terminates in a collection of empty graphs. In the second form (with a plus), it terminates in a collection of complete graphs.
This forms the basis of many algorithms for graph coloring. The ChromaticPolynomial function in the Combinatorica package of the computer algebra system Mathematica uses the second recurrence if the graph is dense, and the first recurrence if the graph is sparse. The worst case running time of either formula satisfies the same recurrence relation as the Fibonacci numbers, so in the worst case, the algorithm runs in time within a polynomial factor of
on a graph with n vertices and m edges. The analysis can be improved to within a polynomial factor of the number of spanning trees of the input graph. In practice, branch and bound strategies and graph isomorphism rejection are employed to avoid some recursive calls, the running time depends on the heuristic used to pick the vertex pair.
Cube method
There is a natural geometric perspective on graph colorings by observing that, as an assignment of natural numbers to each vertex, a graph coloring is a vector in the integer lattice.
Since two vertices and being given the same color is equivalent to the ’th and ’th coordinate in the coloring vector being equal, each edge can be associated with a hyperplane of the form . The collection of such hyperplanes for a given graph is called its graphic arrangement. The proper colorings of a graph are those lattice points which avoid forbidden hyperplanes.
Restricting to a set of colors, the lattice points are contained in the cube . In this context the chromatic polynomial counts the number of lattice points in the -cube that avoid the graphic arrangement.
Computational complexity
The problem of computing the number of 3-colorings of a given graph is a canonical example of a #P-complete problem, so the problem of computing the coefficients of the chromatic polynomial is #P-hard. Similarly, evaluating for given G is #P-complete. On the other hand, for it is easy to compute , so the corresponding problems are polynomial-time computable. For integers the problem is #P-hard, which is established similar to the case . In fact, it is known that is #P-hard for all x (including negative integers and even all complex numbers) except for the three “easy points”. Thus, from the perspective of #P-hardness, the complexity of computing the chromatic polynomial is completely understood.
In the expansion
the coefficient is always equal to 1, and several other properties of the coefficients are known. This raises the question if some of the coefficients are easy to compute. However the computational problem of computing ar for a fixed r ≥ 1 and a given graph G is #P-hard, even for bipartite planar graphs.
No approximation algorithms for computing are known for any x except for the three easy points. At the integer points , the corresponding decision problem of deciding if a given graph can be k-colored is NP-hard. Such problems cannot be approximated to any multiplicative factor by a bounded-error probabilistic algorithm unless NP = RP, because any multiplicative approximation would distinguish the values 0 and 1, effectively solving the decision version in bounded-error probabilistic polynomial time. In particular, under the same assumption, this rules out the possibility of a fully polynomial time randomised approximation scheme (FPRAS). There is no FPRAS for computing for any x > 2, unless NP = RP holds.
| Mathematics | Graph theory | null |
1460629 | https://en.wikipedia.org/wiki/Effective%20temperature | Effective temperature | The effective temperature of a body such as a star or planet is the temperature of a black body that would emit the same total amount of electromagnetic radiation. Effective temperature is often used as an estimate of a body's surface temperature when the body's emissivity curve (as a function of wavelength) is not known.
When the star's or planet's net emissivity in the relevant wavelength band is less than unity (less than that of a black body), the actual temperature of the body will be higher than the effective temperature. The net emissivity may be low due to surface or atmospheric properties, such as the greenhouse effect.
Star
The effective temperature of a star is the temperature of a black body with the same luminosity per surface area () as the star and is defined according to the Stefan–Boltzmann law . Notice that the total (bolometric) luminosity of a star is then , where is the stellar radius. The definition of the stellar radius is obviously not straightforward. More rigorously the effective temperature corresponds to the temperature at the radius that is defined by a certain value of the Rosseland optical depth (usually 1) within the stellar atmosphere. The effective temperature and the bolometric luminosity are the two fundamental physical parameters needed to place a star on the Hertzsprung–Russell diagram. Both effective temperature and bolometric luminosity depend on the chemical composition of a star.
The effective temperature of the Sun is around .
The nominal value defined by the International Astronomical Union for use as a unit of measure of temperature is .
Stars have a decreasing temperature gradient, going from their central core up to the atmosphere. The "core temperature" of the Sun—the temperature at the centre of the Sun where nuclear reactions take place—is estimated to be 15,000,000 K.
The color index of a star indicates its temperature from the very cool—by stellar standards—red M stars that radiate heavily in the infrared to the very hot blue O stars that radiate largely in the ultraviolet. Various colour-effective temperature relations exist in the literature. Their relations also have smaller dependencies on other stellar parameters, such as the stellar metallicity and surface gravity. The effective temperature of a star indicates the amount of heat that the star radiates per unit of surface area. From the hottest surfaces to the coolest is the sequence of stellar classifications known as O, B, A, F, G, K, M.
A red star could be a tiny red dwarf, a star of feeble energy production and a small surface or a bloated giant or even supergiant star such as Antares or Betelgeuse, either of which generates far greater energy but passes it through a surface so large that the star radiates little per unit of surface area. A star near the middle of the spectrum, such as the modest Sun or the giant Capella radiates more energy per unit of surface area than the feeble red dwarf stars or the bloated supergiants, but much less than such a white or blue star as Vega or Rigel.
Planet
Blackbody temperature
To find the effective (blackbody) temperature of a planet, it can be calculated by equating the power received by the planet to the known power emitted by a blackbody of temperature .
Take the case of a planet at a distance from the star, of luminosity .
Assuming the star radiates isotropically and that the planet is a long way from the star, the power absorbed by the planet is given by treating the planet as a disc of radius , which intercepts some of the power which is spread over the surface of a sphere of radius (the distance of the planet from the star). The calculation assumes the planet reflects some of the incoming radiation by incorporating a parameter called the albedo (a). An albedo of 1 means that all the radiation is reflected, an albedo of 0 means all of it is absorbed. The expression for absorbed power is then:
The next assumption we can make is that the entire planet is at the same temperature , and that the planet radiates as a blackbody. The Stefan–Boltzmann law gives an expression for the power radiated by the planet:
Equating these two expressions and rearranging gives an expression for the effective temperature:
Where is the Stefan–Boltzmann constant. Note that the planet's radius has cancelled out of the final expression.
The effective temperature for Jupiter from this calculation is 88 K and 51 Pegasi b (Bellerophon) is 1,258 K. A better estimate of effective temperature for some planets, such as Jupiter, would need to include the internal heating as a power input. The actual temperature depends on albedo and atmosphere effects. The actual temperature from spectroscopic analysis for HD 209458 b (Osiris) is 1,130 K, but the effective temperature is 1,359 K. The internal heating within Jupiter raises the effective temperature to about 152 K.
Surface temperature of a planet
The surface temperature of a planet can be estimated by modifying the effective-temperature calculation to account for emissivity and temperature variation.
The area of the planet that absorbs the power from the star is which is some fraction of the total surface area , where is the radius of the planet. This area intercepts some of the power which is spread over the surface of a sphere of radius . We also allow the planet to reflect some of the incoming radiation by incorporating a parameter called the albedo. An albedo of 1 means that all the radiation is reflected, an albedo of 0 means all of it is absorbed. The expression for absorbed power is then:
The next assumption we can make is that although the entire planet is not at the same temperature, it will radiate as if it had a temperature over an area which is again some fraction of the total area of the planet. There is also a factor , which is the emissivity and represents atmospheric effects. ranges from 1 to 0 with 1 meaning the planet is a perfect blackbody and emits all the incident power. The Stefan–Boltzmann law gives an expression for the power radiated by the planet:
Equating these two expressions and rearranging gives an expression for the surface temperature:
Note the ratio of the two areas. Common assumptions for this ratio are for a rapidly rotating body and for a slowly rotating body, or a tidally locked body on the sunlit side. This ratio would be 1 for the subsolar point, the point on the planet directly below the sun and gives the maximum temperature of the planet — a factor of (1.414) greater than the effective temperature of a rapidly rotating planet.
Also note here that this equation does not take into account any effects from internal heating of the planet, which can arise directly from sources such as radioactive decay and also be produced from frictions resulting from tidal forces.
Earth effective temperature
Earth has an albedo of about 0.306 and a solar irradiance () of at its mean orbital radius of 1.5×108 km. The calculation with ε=1 and remaining physical constants then gives an Earth effective temperature of .
The actual temperature of Earth's surface is an average as of 2020. The difference between the two values is called the greenhouse effect. The greenhouse effect results from materials in the atmosphere (greenhouse gases and clouds) absorbing thermal radiation and reducing emissions to space, i.e., reducing the planet's emissivity of thermal radiation from its surface into space. Substituting the surface temperature into the equation and solving for ε gives an effective emissivity of about 0.61 for a 288 K Earth. Furthermore, these values calculate an outgoing thermal radiation flux of (with ε=0.61 as viewed from space) versus a surface thermal radiation flux of (with ε≈1 at the surface). Both fluxes are near the confidence ranges reported by the IPCC.
| Physical sciences | Basics | Astronomy |
37263123 | https://en.wikipedia.org/wiki/Bag | Bag | A bag (also known regionally as a sack) is a common tool in the form of a non-rigid container, typically made of cloth, leather, bamboo, paper, or plastic. The use of bags predates recorded history, with the earliest bags being lengths of animal skin, cotton, or woven plant fibers, folded up at the edges and secured in that shape with strings of the same material. Bags can be used to carry items such as personal belongings, groceries, and other objects. They come in various shapes and sizes, often equipped with handles or straps for easier carrying.
Bags have been fundamental for the development of human civilization, as they allow people to easily collect and carry loose materials, such as berries or food grains, also allowing them to carry more items in their hands.
The word probably has its origins in the Norse word baggi, from the reconstructed Proto-Indo-European bʰak, but is also comparable to the Welsh baich (load, bundle), and the Greek Τσιαντουλίτσα (Chandulícha, load).
Cheap disposable paper bags and plastic shopping bags are very common, varying in size and strength in the retail trade as a convenience for shoppers, and are often supplied by the shop for free or for a small fee. Customers may also take their own shopping bag(s) to use in shops.
Although paper had been used for wrapping and padding in Ancient China since the 2nd century BC, the first use of paper bags in China (for preserving the flavor of tea) came during the later Tang dynasty (618–907 AD).
History
Bags have been attested for thousands of years and have been used by both men and women. Bags have been prevalent as far back as Ancient Egypt. Many hieroglyphs depict males with bags tied around their waists. The Bible mentions pouches, especially with regard to Judas Iscariot carrying one around, holding his personal items. In the 14th century, wary of pickpockets and thieves, many people used drawstring bags, in which to carry their money. These bags were attached to girdles via a long cord fastened to the waist.
The Australian dillybag is a traditional Australian Aboriginal bag generally woven from plant fibres. Dillybags were and are mainly designed and used by women to gather and transport food, and are most commonly found in the northern parts of Australia.
Women also wore more ornate drawstring bags, typically called hamondeys or tasques, to display their social status. The 14th-century handbags evolved into wedding gifts from groom to bride. These medieval pouches were embroidered, often with depictions of love stories or songs. Eventually, these pouches evolved into what were known as a chaneries, which were used for gaming or food for falcons. During the Renaissance, Elizabethan England's fashions were more ornate than ever before. Women wore their pouches underneath the vast array of petticoats and men wore leather pockets or bagges inside their breeches. Aristocrats began carrying swete bagges filled with sweet-smelling material to make up for poor hygiene.
Modern
In the modern world, bags are ubiquitous, with many people routinely carrying a wide variety of them in the form of cloth or leather briefcases, handbags, and backpacks, and with bags made from more disposable materials such as paper or plastic being used for shopping or to carry groceries. Today, bags are also used as a fashion statement. A bag may be closable by a zipper, snap fastener, etc., or simply by folding (e.g. in the case of a paper bag). Sometimes a money bag or travel bags has a lock. The bag likely predates its inflexible variant, the basket, and usually has the additional advantage of being foldable or otherwise compressible to smaller sizes. On the other hand, baskets, being made of more rigid materials, may be better at protecting their contents.
An empty bag may or may not be very light and foldable to a small size. If it is, this is convenient for carrying it to the place where it is needed, such as a shop, and for storage of empty bags. Bags vary from small ones, like purses, to larger ones used for traveling such as a suitcase. The pockets of clothing are also a kind of bag, built into the clothing for the carrying of suitably small objects.
Environmental aspects
There are environmental concerns regarding use and disposal of plastic bags. Efforts are being taken to control and reduce their use in some European Union countries, including Ireland and the Netherlands. In some cases these cheap bags are taxed so the customer must pay a fee where they may not have done previously. Sometimes heavy duty reusable plastic and fabric bags are sold, typically costing €0.50 to €1, and these may replace disposable bags entirely. Sometimes free replacements are offered when the bag wears out. The UK has charged 5p per plastic carrier bag in larger shops since 2015. This trend has spread to some cities in the United States. Recently many countries have banned the use of plastic bags. Paper bags emerge as a great replacement for plastic bags; however, paper bags tend to be more expensive.
A bag may or may not be disposable; however, even a disposable bag can often be used many times, for economic and environmental reasons. On the other hand, there may be logistic or hygienic reasons to use a bag only once. For example, a garbage bag is often disposed of with its contents. A bag for packaging a disposable product is often disposed of when empty. Similarly, bags used as receptacles in medical procedures, such as the colostomy bag used to collect waste from a surgically diverted biological system, are typically disposed of as medical waste. Many snack foods, such as pretzels, cookies, and potato chips, are available in disposable single-use sealed bags.
Types of bags
Antistatic bag (used for shipping electronic components)
Backpack
Bag-in-box
Baguette
Bin bag, garbage bag, or trash bag
Blue bag
Biodegradable bag
Bivouac bag
Body bag
Book bag
Booster bag
Bota bag
Bulk bag, another name for a flexible intermediate bulk container
Burn bag
Camera bag
Carpet bag
Cooler bag
Diaper bag
Diplomatic bag
Douche bag
Duffel bag
Dunnage bag
Flour sack
Gaji bag
Garment bag
Gladstone bag
Gunny sack
Handbag or purse
Hobo bag
Ita-bag
Lifting bag
Mail bag
Messenger bag
Millbank bag
Money bag
Paper bag
Plastic bag
Popcorn bag
Punching bag
Sandbag
Satchel
Security bag
Sling bag, worn over the shoulder
Shopping bag
Stuff sack
Suicide bag
Thermal bag
Tote bag
Travel bag or suitcase
Tucker bag
Other
Airbag (vehicle safety device)
Bagpipes
Bag valve mask
Bean bag
Bag valve mask
Bota bag
Bulgarian Bag
Coin purse
Coffee bag
Ita-bag
Milk bag
Oven bag
Pastry bag
Punching bag (a piece of physical training equipment)
Perhaps-bag or Netted sack
Portable hyperbaric bag
Raschen bag
Sachet
Sleeping bag
Sonali Bag
Spice bag
Tea bag
Vacuum bag
Zipper storage bag
| Technology | Containers | null |
2927640 | https://en.wikipedia.org/wiki/Acute%20medicine | Acute medicine | Acute medicine, also known as acute internal medicine (AIM), is a specialty within internal medicine concerned with the immediate and early specialist management of adult patients with a wide range of medical conditions who present in hospital as emergencies. It developed in the United Kingdom in the early 2000s as a dedicated field of medicine, together with the establishment of acute medical units in numerous hospitals. Acute medicine is distinct from the broader field of emergency medicine, which is concerned with the management of all people attending the emergency department, not just those with internal medicine diagnoses.
History
The field developed in the United Kingdom after the Royal College of Physicians of Edinburgh and the Royal College of Physicians and Surgeons of Glasgow published a joint report in 1998 emphasising the importance of appropriate care for people with acute medical problems. Further reports led to the development of acute medicine as a dedicated specialty, and in 2003 it was recognised by the Specialist Training Authority as a subspecialty of General Internal Medicine.
Around the same time, it was recognised that care for acutely admitted patients should ideally be concentrated in "medical assessment units" (MAUs), later named "acute medical units" (AMUs). A physician experienced in the management of acute medical problems could assess and treat these patients in the most appropriate fashion for the first 48 hours of their admission, aiming either for an early discharge with appropriate outpatient follow-up or transfer to a specialist ward. Severely ill patients who need close observation but do not require intensive care may be treated in a dedicated area such as a physician-run high-dependency unit.
In 2007, some questioned whether the specialty would have a long-term future, if at some point UK government ED targets ceased to exist. However, Robert Wachter has stated that acute medical units "have been associated with lower inpatient mortality, improved patient and staff satisfaction, reduced hospital stays, and increased throughput."
A further development has been the increase of ambulatory care. Where patients were previously admitted to hospital, it may now be possible for them to attend a clinic or an assessment area a number of times while their progress is monitored. This is now a very common approach to suspected deep vein thrombosis, but the NHS Institute for Innovation and Improvement has identified a number of other conditions that can be managed in an ambulatory emergency care setting.
In 2009, the General Medical Council approved acute medicine as a distinct specialty, allowing doctors to specialise in it in order to receive their Certificate of Completion of Training.
Organisations
The Society for Acute Medicine was founded in 2000. It is "the national representative body for staff caring for medical patients in the acute hospital setting. In the Netherlands, the Dutch Acute Medicine (DAM) society was formed in 2012 and held its first Congress on 28 September 2012 in the VU University Medical Center in Amsterdam. Its curriculum aimed to ensure 12 training posts throughout the Netherlands.
| Biology and health sciences | Fields of medicine | Health |
2930244 | https://en.wikipedia.org/wiki/Atlantic%20hurricane%20season | Atlantic hurricane season | The Atlantic hurricane season is the period in a year, from June 1 through November 30, when tropical or subtropical cyclones are most likely to form in the North Atlantic Ocean. These dates, adopted by convention, encompass the period in each year when most tropical cyclogenesis occurs in the basin. Even so, subtropical or tropical cyclogenesis is possible at any time of the year, and often does occur.
Worldwide, a season's climatological peak activity takes place in late summer, when the difference between air temperature and sea surface temperatures is the greatest. Peak activity in an Atlantic hurricane season happens from late August through September, with a midpoint on September 10.
Atlantic tropical and subtropical cyclones that reach tropical storm intensity are named from a predetermined list. On average, 14 named storms occur each season, with an average of 7 becoming hurricanes and 3 becoming major hurricanes, Category 3 or higher on the Saffir–Simpson scale. The most active season on record was 2020, during which 30 named tropical cyclones formed. Despite this, the 2005 season had more hurricanes, developing a record of 15 such storms. The least active season was 1914, with only one known tropical cyclone developing during that year.
Concept
The understanding that Atlantic hurricanes are most commonplace during a certain period of the year has been long recognized. Historical delineations of the Atlantic hurricane season varied but generally covered some part of the estival (summer) and autumnal months. Some early descriptions of the season's bounds theorized that the timing of the full moon or the moon's phases as a whole could be used to more precisely delineate the hurricane season. In the second volume of Voyages and Descriptions (published in 1700), English explorer and naturalist William Dampier observed that hurricanes in the Caribbean Sea were expected in July, August, and September. Mariners in the 18th century generally regarded the period from July to the end of October as the "hurricane season" based on the frequency of storms striking the Caribbean islands and the trajectories of ships traversing the Atlantic.
The hurricane season was also an important influence on European naval operations within the West Indies, forcing the movement of materiel to be expedited before its onset or delayed until its end. English admiral Edward Vernon described the "hurricane months" of August and September within the West Indies as a particularly vulnerable time for maritime logistics; Vernon argued that the most optimal time for a fleet to be dispatched from Great Britain to attack Spanish assets in the Americas was August or September, in part because such ships would more likely avoid hurricanes by the time they reached the West Indies. American geographer Jedidiah Morse defined the hurricane season as the months of August, September, and October in his treatise The American Universal Geography. American meteorologist William Charles Redfield defined the hurricane season as lasting from July 15 to October 15, citing the timeframe during which some insurance underwriters raised premiums in response to the increased likelihood of hurricanes. Based on a catalog of 355 storms between 1493–1855 in the North Atlantic compiled by M. André Poëy, W. H. Rosser described the months of July, August, September, and October as comprising the "true hurricane season of the West Indies" in his 1876 book The Law of Storms Considered Practically.
The concept of the hurricane season took on a more practical significance in forecasting operations as the United States Weather Bureau began to extend its weather prediction efforts and data collection into the tropics. In 1882, the bureau briefly considered an effort to adopt special hurricane signals between July and October 20 to emphasize the danger of such storms during that period, but dropped the effort due to a lack of funding. When the U.S. Weather Bureau built a network of weather observatories in the Caribbean in 1898, these sites telegraphed weather observations at 8 a.m. daily to the bureau's regional headquarters in Havana, Cuba, during the hurricane season; this season was defined as lasting from the beginning of June through October. By 1907, these stations in the West Indies operated within a hurricane season defined as beginning on June 15 and ending on November 15. The starting date of these regular reports was moved back to June 1 by 1915. In 1917, an increase in funding for the U.S. Weather Bureau's observing networks in the Caribbean region led to these stations reporting twice daily during a hurricane season expanded to cover the June 1 to November 30 period. This delineation was maintained when the bureau (in cooperation with United Fruit Company) began to broadcast special weather bulletins for Caribbean shipping during the hurricane season in 1922, providing information on active hurricanes and warnings twice daily.
The basic concept of an official hurricane season began during 1935, when dedicated wire circuits known as hurricane circuits began to be set up along the Gulf and Atlantic coasts, a process completed by 1955. It was originally the time frame when the tropics were monitored routinely for tropical cyclone activity, and was originally defined as from June 15 through October 31. Over the years, the beginning date was shifted back to June 1, while the end date was shifted to November 15, before settling at November 30 by 1965. This was when hurricane reconnaissance planes were sent out to fly across the Atlantic and Gulf of Mexico on a routine basis to look for potential tropical cyclones, in the years before the continuous weather satellite era. Since regular satellite surveillance began, hurricane hunter aircraft fly only into storm areas which are first spotted by satellite imagery. The six-month official hurricane season established in 1965 by the National Hurricane Center (NHC) remains the current delineation of the Atlantic hurricane season. These bounds contain over 97 percent of Atlantic tropical cyclone activity. While this definition was chosen in part to make it easier for the public to remember the timing of hurricanes, storms have often formed outside the official seasonal bounds. Following several consecutive years of Atlantic tropical cyclones developing before the official June 1 start date, the World Meteorological Organization recommended in 2021 that the NHC assess moving the start date to May 15. In response, the NHC formed a team to develop quantiative criteria to evaluate extending the seasonal bounds. The agency's routine tropical weather outlooks, historically issued during the hurricane season beginning on June 1, were instead started on May 15 beginning in 2021.
Operations
During the hurricane season, the National Hurricane Center routinely issues their Tropical Weather Outlook product, which identifies areas of concern within the tropics which could develop into tropical cyclones. If systems occur outside the defined hurricane season, special Tropical Weather Outlooks will be issued. Routine coordination occurs at 1700 UTC each day between the Weather Prediction Center and National Hurricane Center to identify systems for the pressure maps three to seven days into the future within the tropics, and points for existing tropical cyclones six to seven days into the future. Possible tropical cyclones are depicted with a closed isobar, while systems with less certainty to develop are depicted as "spot lows" with no isobar surrounding them.
HURDAT
The North Atlantic hurricane database, or HURDAT, is the database for all tropical storms and hurricanes for the Atlantic Ocean, Gulf of Mexico and Caribbean Sea, including those that have made landfall in the United States. The original database of six-hourly positions and intensities was put together in the 1960s in support of the Apollo space program to help provide statistical track forecast guidance. In the intervening years, this databasewhich is now freely and easily accessible on the Internet from the National Hurricane Center's (NHC) webpagehas been utilized for a wide variety of uses: climatic change studies, seasonal forecasting, risk assessment for county emergency managers, analysis of potential losses for insurance and business interests, intensity forecasting techniques and verification of official and various model predictions of track and intensity.
HURDAT was not designed with all of these uses in mind when it was first put together and not all of them may be appropriate given its original motivation. HURDAT contains numerous systematic as well as some random errors in the database. Additionally, analysis techniques have changed over the years at NHC as their understanding of tropical cyclones has developed, leading to biases in the historical database. Another difficulty in applying the hurricane database to studies concerned with landfalling events is the lack of exact location, time and intensity at hurricane landfall.
Re-analysis project
HURDAT is regularly updated annually to reflect the previous season's activity. The older portion of the database has been regularly revised since 2001. The first time in 2001 led to the addition of tropical cyclone tracks for the years 1851 to 1885. The second time was in October 2002 when Hurricane Andrew (August 1992) was upgraded to a Category 5. Recent efforts into uncovering undocumented historical hurricanes in the late 19th and 20th centuries by various researchers have greatly increased our knowledge of these past events. Tropical storms from 1851 to 1970 have already been reanalyzed with most recently, re-analysis of tropical storms from 1961 to 1965 being completed and integrated into HURDAT database in November 2019, and re-analysis of tropical storms from 1966 to 1970 being completed and integrated into HURDAT database in January 2022. Possible changes for the years 1971 onward are not yet incorporated into the HURDAT database. Due to these issues, a re-analysis of the Atlantic hurricane database is being attempted that will be completed in three years.
In addition to the groundbreaking work by Partagas Cigars, additional analyses, digitization and quality control of the data was carried out by researchers at the NOAA Hurricane Research Division funded by the NOAA Office of Global Programs.
The National Hurricane Center's Best Track Change Committee has approved changes for a few recent cyclones, such as Hurricane Andrew. Official changes to the Atlantic hurricane database are approved by the National Hurricane Center Best Track Change Committee.
1494–1850 (pre-HURDAT era)
1851–1899 (within HURDAT data)
1850s
1860s
1870s
1880s
1890s
1900s
NOTE: In the following tables, all estimates of damage costs are expressed in contemporaneous US dollars (USD).
1900s
1910s
1920s
1930s
1940s
1950s
1960s
1970s
1980s
1990s
2000s
NOTE: In the following tables, all estimates of damage costs are expressed in contemporaneous US dollars (USD).
2000s
2010s
2020s
Climatology
A 2011 study analyzing one of the main sources of hurricanesthe African easterly wave (AEW)found that the change in AEWs is closely linked to increased activity of intense hurricanes in the North Atlantic. The synoptic concurrence of AEWs in driving the dynamics of the Sahel greening also appears to increase tropical cyclogenesis over the North Atlantic.
| Physical sciences | Seasons | Earth science |
2930734 | https://en.wikipedia.org/wiki/Distributive%20shock | Distributive shock | Distributive shock is a medical condition in which abnormal distribution of blood flow in the smallest blood vessels results in inadequate supply of blood to the body's tissues and organs. It is one of four categories of shock, a condition where there is not enough oxygen-carrying blood to meet the metabolic needs of the cells which make up the body's tissues and organs. Distributive shock is different from the other three categories of shock in that it occurs even though the output of the heart is at or above a normal level. The most common cause is sepsis leading to a type of distributive shock called septic shock, a condition that can be fatal.
Types
Elbers and Ince have identified five classes of abnormal microcirculatory flow in distributive shock using side stream dark field microscopy.
Class I: all capillaries are stagnant when there is normal or sluggish venular flow.
Class II: there are empty capillaries next to capillaries that have flowing red blood cells.
Class III: there are stagnant capillaries next to capillaries with normal blood flow.
Class IV: hyperdynamic flow in capillaries adjacent to capillaries that are stagnant.
Class V: widespread hyperdynamic flow in the microcirculatory system.
According to the cause, there are 4 types of distributive shock:
Neurogenic shock: Decreased sympathetic stimulation leading to decreased vessel tone.
Anaphylactic shock
Septic shock
Shock due to adrenal crisis
Causes
In addition to sepsis, distributive shock can be caused by systemic inflammatory response syndrome (SIRS) due to conditions other than infection such as pancreatitis, burns or trauma. Other causes include toxic shock syndrome (TSS), anaphylaxis (a sudden, severe allergic reaction), adrenal insufficiency, reactions to drugs or toxins, heavy metal poisoning, hepatic (liver) insufficiency and damage to the central nervous system. Causes of adrenal insufficiency leading to distributive shock include acute worsening of chronic adrenal insufficiency, destruction or removal of the adrenal glands, suppression of adrenal gland function due to exogenous steroids, hypopituitarism and metabolic failure of hormone production.
Pathophysiology
The cause of inadequate tissue perfusion (blood delivery to tissues) in distributive shock is a lack of normal responsiveness of blood vessels to vasoconstrictive agents and direct vasodilation.
There are four types of distributive shock. The most common, septic shock, is caused by an infection, most frequently by bacteria, but viruses, fungi and parasites have been implicated. Infection sites most likely to lead to septic shock are chest, abdomen and genitourinary tract. In septic shock the blood flow in the microvasculature is abnormal with some capillaries underperfused and others with normal to high blood flow. The endothelial cells lining the blood vessels become less responsive to vasoconstrictive agents, lose their glycocalyx (normal coating) and negative ionic charge, become leaky and cause extensive over-expression of nitric oxide. The coagulation cascade is also disrupted. Tissue factor that initiates the clotting cascade is produced by activated monocytes and the endothelial cells lining the blood vessels while antithrombin and fibrinolysis are impaired. Disseminated intravascular coagulation (DIC) can result from the thrombin produced in the inflammatory response. The ability of red blood cells to change shape decreases and their tendency to clump together increases, inhibiting their flow through the microvasculature.
In anaphylactic shock low blood pressure is related to decreased systemic vascular resistance (SVR) triggered primarily by a massive release of histamine by mast cells activated by antigen-bound immunoglobulin E and also by increased production and release of prostaglandins.
Neurogenic shock is caused by the loss of vascular tone normally supported by the sympathetic nervous system due to injury to the central nervous system especially spinal cord injury. Rupture of a hollow organ, with subsequent evacuation of contents in the peritoneal cavity could also determine neurogenic shock, a subtype of distributive shock. This happens due to the widespread peritoneal irritation by the ruptured viscus contents, as in peptic ulcer perforation, with consequent strong vagal activation, and generalized, extensive peripheral vasodilation and bradycardia. Thus, in neurogenic shock, there is decreased systemic vascular resistance, CVP is typically decreased, CO decreased or normal, and PAOP decreased.
Distributive shock associated with adrenal crisis results from inadequate steroid hormones.
Diagnosis
Treatment
The main goals of treatment in distributive shock are to reverse the underlying cause and achieve hemodynamic stabilization. Immediate treatment involves fluid resuscitation and the use of vasoactive drugs, both vasopressors and inotropes. Hydrocortisone is used for people whose hypotension does not respond to fluid resuscitation and vasopressors. Opening and keeping open the microcirculation is a consideration in the treatment of distributive shock, as a result limiting the use of vasopressors has been suggested. Control of inflammation, vascular function and coagulation to correct pathological differences in blood flow and microvascular shunting has been pointed to as a potentially important adjunct goal in the treatment of distributive shock.
People with septic shock are treated with antimicrobial drugs to treat the causative infection. Some sources of infection require surgical intervention including necrotizing fasciitis, cholangitis, abscess, intestinal ischemia, or infected medical devices.
Anaphylactic shock is treated with epinephrine. The use of vasopressin together with norepinephrine rather than norepinephrine alone appears to decrease the risk of atrial fibrillation but with few other benefits.
Prognosis
Septic shock is associated with significant mortality and is the leading non-cardiac cause of death in intensive care units (ICUs).
Research directions
The choice of fluids for resuscitation remains an area of research, the Surviving Sepsis Campaign an international consortium of experts, did not find adequate evidence to support the superiority of crystalloid fluids versus colloid fluids. Drugs such as pyridoxalated hemoglobin polyoxyethylene, which scavenge nitric oxide from the blood, have been investigated, as well as methylene blue which may inhibit the nitric oxide-cyclic guanosine monophosphate (NO-cGMP) pathway which has been suggested to play a significant role in distributive shock.
| Biology and health sciences | Cardiovascular disease | Health |
35888611 | https://en.wikipedia.org/wiki/Semibatch%20reactor | Semibatch reactor | For both chemical and biological engineering, Semibatch (semiflow) reactors operate much like batch reactors in that they take place in a single stirred tank with similar equipment. However, they are modified to allow reactant addition and/or product removal in time.
A normal batch reactor is filled with reactants in a single stirred tank at time and the reaction proceeds. A semi batch reactor, however, allows partial filling of reactants with the flexibility of adding more as time progresses. Stirring in both types is very efficient, which allows batch and semi batch reactors to assume a uniform composition and temperature throughout.
Advantages
The flexibility of adding more reactants over time through semi batch operation has several advantages over a batch reactor. These include:
Improved selectivity of a reaction
Sometimes a particular reactant can go through parallel paths that yield two different products, only one of which is desired. Consider the simple example below:
A→U (desired product)
A→W (undesired product)
The rate expressions, considering the variability of the volume of reaction, are:
= -
= -
= -
Where is the molar rate of addition of the reactant A. Note that the presence of these addition terms, which could be negative in case of products removal (e.g. by fractional distillation) are the ones marking the difference of the semi batch reactor cases from the simpler batch cases.
For standard batch reactors (no addition terms) the selectivity of the desired product is defined as:
S = =
S = = for constant volume (i.e. batch) reactions.
If , the concentration of the reactant should be kept at a low level in order to maximize selectivity. This can be accomplished using a semibatch reactor.
Better control of exothermic reactions
Exothermic reactions release heat, and ones that are highly exothermic can cause safety concerns. Semibatch reactors allow for slow addition of reactants in order to control the heat released and thus, temperature, in the reactor.
Product removal through a purge stream
In order to minimize the reversibility of a reaction one must minimize the concentration of the product. This can be done in a semibatch reactor by using a purge stream to remove products and increase the net reaction rate by favoring the forward reaction.
Reactor choice
It is important to understand that these advantages are more applicable to the decision between using a batch, a semibatch or a continuous reactor in a certain process. Both batch and semibatch reactors are more suitable for liquid phase reactions and small scale production, because they usually require lower capital costs than a continuously stirred tank reactor operation (CSTR), but incur greater costs per unit if production needs to be scaled up. These per unit costs include labor, materials handling (filling, emptying, cleaning), protective measures, and nonproductive periods that result from changeovers when switching batches. Hence, the capital costs must be weighed against operating costs to determine the correct reactor design to be implemented.
| Physical sciences | Chemical engineering | Chemistry |
35895879 | https://en.wikipedia.org/wiki/Effects%20of%20climate%20change%20on%20oceans | Effects of climate change on oceans | There are many effects of climate change on oceans. One of the most important is an increase in ocean temperatures. More frequent marine heatwaves are linked to this. The rising temperature contributes to a rise in sea levels due to the expansion of water as it warms and the melting of ice sheets on land. Other effects on oceans include sea ice decline, reducing pH values and oxygen levels, as well as increased ocean stratification. All this can lead to changes of ocean currents, for example a weakening of the Atlantic meridional overturning circulation (AMOC). The main cause of these changes are the emissions of greenhouse gases from human activities, mainly burning of fossil fuels and deforestation. Carbon dioxide and methane are examples of greenhouse gases. The additional greenhouse effect leads to ocean warming because the ocean takes up most of the additional heat in the climate system. The ocean also absorbs some of the extra carbon dioxide that is in the atmosphere. This causes the pH value of the seawater to drop. Scientists estimate that the ocean absorbs about 25% of all human-caused emissions.
The various layers of the oceans have different temperatures. For example, the water is colder towards the bottom of the ocean. This temperature stratification will increase as the ocean surface warms due to rising air temperatures. Connected to this is a decline in mixing of the ocean layers, so that warm water stabilises near the surface. A reduction of cold, deep water circulation follows. The reduced vertical mixing makes it harder for the ocean to absorb heat. So a larger share of future warming goes into the atmosphere and land. One result is an increase in the amount of energy available for tropical cyclones and other storms. Another result is a decrease in nutrients for fish in the upper ocean layers. These changes also reduce the ocean's capacity to store carbon. At the same time, contrasts in salinity are increasing. Salty areas are becoming saltier and fresher areas less salty.
Warmer water cannot contain the same amount of oxygen as cold water. As a result, oxygen from the oceans moves to the atmosphere. Increased thermal stratification may reduce the supply of oxygen from surface waters to deeper waters. This lowers the water's oxygen content even more. The ocean has already lost oxygen throughout its water column. Oxygen minimum zones are increasing in size worldwide.
These changes harm marine ecosystems, and this can lead to biodiversity loss or changes in species distribution. This in turn can affect fishing and coastal tourism. For example, rising water temperatures are harming tropical coral reefs. The direct effect is coral bleaching on these reefs, because they are sensitive to even minor temperature changes. So a small increase in water temperature could have a significant impact in these environments. Another example is loss of sea ice habitats due to warming. This will have severe impacts on polar bears and other animals that rely on it. The effects of climate change on oceans put additional pressures on ocean ecosystems which are already under pressure by other impacts from human activities.
Changes due to rising greenhouse gas levels
Presently (2020), atmospheric carbon dioxide (CO2) levels of more than 410 parts per million (ppm) are nearly 50% higher than preindustrial levels. These elevated levels and rapid growth rates are unprecedented in the geological record's 55 million years. The source for this excess is clearly established as human-driven, reflecting a mix of fossil fuel burning, industrial, and land-use/land-change emissions. The idea that the ocean serves as a major sink for anthropogenic has been discussed in scientific literature since at least the late 1950s. Several pieces of evidence point to the ocean absorbing roughly a quarter of total anthropogenic emissions.
The latest key findings about the observed changes and impacts from 2019 include:
Rising ocean temperature
It is clear that the ocean is warming as a result of climate change, and this rate of warming is increasing. The global ocean was the warmest it had ever been recorded by humans in 2022. This is determined by the ocean heat content, which exceeded the previous 2021 maximum in 2022. The steady rise in ocean temperatures is an unavoidable result of the Earth's energy imbalance, which is primarily caused by rising levels of greenhouse gases. Between pre-industrial times and the 2011–2020 decade, the ocean's surface has heated between 0.68 and 1.01 °C.
The majority of ocean heat gain occurs in the Southern Ocean. For example, between the 1950s and the 1980s, the temperature of the Antarctic Southern Ocean rose by 0.17 °C (0.31 °F), nearly twice the rate of the global ocean.
The warming rate varies with depth. The upper ocean (above 700 m) is warming the fastest. At an ocean depth of a thousand metres the warming occurs at a rate of nearly 0.4 °C per century (data from 1981 to 2019). In deeper zones of the ocean (globally speaking), at 2000 metres depth, the warming has been around 0.1 °C per century. The warming pattern is different for the Antarctic Ocean (at 55°S), where the highest warming (0.3 °C per century) has been observed at a depth of 4500 m.
Marine heatwaves
Marine heatwaves also take their toll on marine life: For example, due to fall-out from the 2019-2021 Pacific Northwest marine heatwave, Bering Sea snow crab populations declined 84% between 2018 and 2022, a loss of 9.8 billion crabs.
Ocean heat content
The ocean temperature varies from place to place. Temperatures are higher near the equator and lower at the poles. As a result, changes in total ocean heat content best illustrate ocean warming. When compared to 1969–1993, heat uptake has increased between 1993 and 2017.
Ocean acidification
Time scales
Many ocean-related elements of the climate system respond slowly to warming. For instance, acidification of the deep ocean will continue for millennia, and the same is true for the increase in ocean heat content. Similarly, sea level rise will continue for centuries or even millennia even if greenhouse gas emissions are brought to zero, due to the slow response of ice sheets to warming and the continued uptake of heat by the oceans, which expand when warmed.
Effects on the physical environment
Sea level rise
Many coastal cities will experience coastal flooding in the coming decades and beyond. Local subsidence, which may be natural but can be increased by human activity, can exacerbate coastal flooding. Coastal flooding will threaten hundreds of millions of people by 2050, particularly in Southeast Asia.
Changing ocean currents
Ocean currents are caused by temperature variations caused by sunlight and air temperatures at various latitudes, as well as prevailing winds and the different densities of salt and fresh water. Warm air rises near the equator. Later, as it moves toward the poles, it cools again. Cool air sinks near the poles, but warms and rises again as it moves toward the equator. This produces Hadley cells, which are large-scale wind patterns, with similar effects driving a mid-latitude cell in each hemisphere. Wind patterns associated with these circulation cells drive surface currents which push the surface water to higher latitudes where the air is colder. This cools the water, causing it to become very dense in comparison to lower latitude waters, causing it to sink to the ocean floor, forming North Atlantic Deep Water (NADW) in the north and Antarctic Bottom Water (AABW) in the south.
Driven by this sinking and the upwelling that occurs in lower latitudes, as well as the driving force of the winds on surface water, the ocean currents act to circulate water throughout the sea. When global warming is factored in, changes occur, particularly in areas where deep water is formed. As the oceans warm and glaciers and polar ice caps melt, more and more fresh water is released into the high latitude regions where deep water forms, lowering the density of the surface water. As a result, the water sinks more slowly than it would normally.
The Atlantic Meridional Overturning Circulation (AMOC) may have weakened since the preindustrial era, according to modern observations and paleoclimate reconstructions (the AMOC is part of a global thermohaline circulation), but there is too much uncertainty in the data to know for certain. Climate change projections assessed in 2021 indicate that the AMOC is very likely to weaken over the course of the 21st century. A weakening of this magnitude could have a significant impact on global climate, with the North Atlantic being particularly vulnerable.
Any changes in ocean currents affect the ocean's ability to absorb carbon dioxide (which is affected by water temperature) as well as ocean productivity because the currents transport nutrients (see Impacts on phytoplankton and net primary production). Because the AMOC deep ocean circulation is slow (it takes hundreds to thousands of years to circulate the entire ocean), it is slow to respond to climate change.
Increasing stratification
Changes in ocean stratification are significant because they can influence productivity and oxygen levels. The separation of water into layers based on density is known as stratification. Stratification by layers occurs in all ocean basins. The stratified layers limit how much vertical water mixing takes place, reducing the exchange of heat, carbon, oxygen and particles between the upper ocean and the interior. Since 1970, there has been an increase in stratification in the upper ocean due to global warming and, in some areas, salinity changes. The salinity changes are caused by evaporation in tropical waters, which results in higher salinity and density levels. Meanwhile, melting ice can cause a decrease in salinity at higher latitudes.
Temperature, salinity and pressure all influence water density. As surface waters are often warmer than deep waters, they are less dense, resulting in stratification. This stratification is crucial not just in the production of the Atlantic Meridional Overturning Circulation, which has worldwide weather and climate ramifications, but it is also significant because stratification controls the movement of nutrients from deep water to the surface. This increases ocean productivity and is associated with the compensatory downward flow of water that carries oxygen from the atmosphere and surface waters into the deep sea.
Reduced oxygen levels
Climate change has an impact on ocean oxygen, both in coastal areas and in the open ocean.
The open ocean naturally has some areas of low oxygen, known as oxygen minimum zones. These areas are isolated from the atmospheric oxygen by sluggish ocean circulation. At the same time, oxygen is consumed when sinking organic matter from surface waters is broken down. These low oxygen ocean areas are expanding as a result of ocean warming which both reduces water circulation and also reduces the oxygen content of that water, while the solubility of oxygen declines as the temperature rises.
Overall ocean oxygen concentrations are estimated to have declined 2% over 50 years from the 1960s. The nature of the ocean circulation means that in general these low oxygen regions are more pronounced in the Pacific Ocean. Low oxygen represents a stress for almost all marine animals. Very low oxygen levels create regions with much reduced fauna. It is predicted that these low oxygen zones will expand in future due to climate change, and this represents a serious threat to marine life in these oxygen minimum zones.
The second area of concern relates to coastal waters where increasing nutrient supply from rivers to coastal areas leads to increasing production and sinking organic matter which in some coastal regions leads to extreme oxygen depletion, sometimes referred to as dead zones. These dead zones are expanding driven particularly by increasing nutrient inputs, but also compounded by increasing ocean stratification driven by climate change.
Oceans turning green
Satellite image analysis reveals that the oceans have been gradually turning green from blue as climate breakdown continues. The color change has been detected for a majority of the word's ocean surfaces and may be due to changing plankton populations caused by climate change.
Changes to Earth's weather system and wind patterns
Climate change and the associated warming of the ocean will lead to widespread changes to the Earth's climate and weather system including increased tropical cyclone and monsoon intensities and weather extremes with some areas becoming wetter and others drier. Changing wind patterns are predicted to increase wave heights in some areas.
Intensifying tropical cyclones
Human-induced climate change "continues to warm the oceans which provide the memory of past accumulated effects". The result is a higher ocean heat content and higher sea surface temperatures. In turn, this "invigorates tropical cyclones to make them more intense, bigger, longer lasting and greatly increases their flooding rains". One example is Hurricane Harvey in 2017.
Salinity changes
Due to global warming and increased glacier melt, thermohaline circulation patterns may be altered by increasing amounts of freshwater released into oceans and, therefore, changing ocean salinity. Thermohaline circulation is responsible for bringing up cold, nutrient-rich water from the depths of the ocean, a process known as upwelling.
Seawater consists of fresh water and salt, and the concentration of salt in seawater is called salinity. Salt does not evaporate, thus the precipitation and evaporation of freshwater influences salinity strongly. Changes in the water cycle are therefore strongly visible in surface salinity measurements, which has been known since the 1930s.
The long term observation records show a clear trend: the global salinity patterns are amplifying in this period. This means that the high saline regions have become more saline, and regions of low salinity have become less saline. The regions of high salinity are dominated by evaporation, and the increase in salinity shows that evaporation is increasing even more. The same goes for regions of low salinity that are becoming less saline, which indicates that precipitation is becoming more intensified.
Sea ice decline and changes
Sea ice decline occurs more in the Arctic than in Antarctica, where it is more a matter of changing sea ice conditions.
Impacts on biological processes
Ocean productivity
The process of photosynthesis in the surface ocean releases oxygen and consumes carbon dioxide. This photosynthesis in the ocean is dominated by phytoplankton – microscopic free-floating algae. After the plants grow, bacterial decomposition of the organic matter formed by photosynthesis in the ocean consumes oxygen and releases carbon dioxide. The sinking and bacterial decomposition of some organic matter in deep ocean water, at depths where the waters are out of contact with the atmosphere, leads to a reduction in oxygen concentrations and increase in carbon dioxide, carbonate and bicarbonate. This cycling of carbon dioxide in oceans is an important part of the global carbon cycle.
The photosynthesis in surface waters consumes nutrients (e.g. nitrogen and phosphorus) and transfers these nutrients to deep water as the organic matter produced by photosynthesis sinks upon the death of the organisms. Productivity in surface waters therefore depends in part on the transfer of nutrients from deep water back to the surface by ocean mixing and currents. The increasing stratification of the oceans due to climate change therefore acts generally to reduce ocean productivity. However, in some areas, such as previously ice covered regions, productivity may increase. This trend is already observable and is projected to continue under current projected climate change. In the Indian Ocean for example, productivity is estimated to have declined over the past sixty years due to climate warming and is projected to continue.
Ocean productivity under a very high emission scenario (RCP8.5) is very likely to drop by 4-11% by 2100. The decline will show regional variations. For example, the tropical ocean NPP will decline more: by 7–16% for the same emissions scenario. Less organic matter will likely sink from the upper oceans into deeper ocean layers due to increased ocean stratification and a reduction in nutrient supply. The reduction in ocean productivity is due to the "combined effects of warming, stratification, light, nutrients and predation".
Calcifying organisms and ocean acidification
Harmful algal blooms
Although the drivers of harmful algal blooms (HABs) are poorly understood, they appear to have increased in range and frequency in coastal areas since the 1980s. This is the result of human induced factors such as increased nutrient inputs (nutrient pollution) and climate change (in particular the warming of water temperatures). The parameters that affect the formation of HABs are ocean warming, marine heatwaves, oxygen loss, eutrophication and water pollution. These increases in HABs are of concern because of the impact of their occurrence on local food security, tourism and the economy.
It is however also possible that the perceived increase in HABs globally is simply due to more severe bloom impacts and better monitoring and not due to climate change.
Impacts on coral reefs and fisheries
Coral reefs
While some mobile marine species can migrate in response to climate change, others such as corals find this much more difficult. A coral reef is an underwater ecosystem characterised by reef-building corals. Reefs are formed by colonies of coral polyps held together by calcium carbonate. Coral reefs are important centres of biodiversity and vital to millions of people who rely on them for coastal protection, food and for sustaining tourism in many regions.
Warm water corals are clearly in decline, with losses of 50% over the last 30–50 years due to multiple threats from ocean warming, ocean acidification, pollution and physical damage from activities such as fishing. These pressures are expected to intensify.
The warming ocean surface waters can lead to bleaching of the corals which can cause serious damage and/or coral death. The IPCC Sixth Assessment Report in 2022 found that: "Since the early 1980s, the frequency and severity of mass coral bleaching events have increased sharply worldwide". Marine heatwaves have caused coral reef mass mortality. It is expected that many coral reefs will suffer irreversible changes and loss due to marine heatwaves with global temperatures increasing by more than 1.5 °C.
Coral bleaching occurs when thermal stress from a warming ocean results in the expulsion of the symbiotic algae that resides within coral tissues. These symbiotic algae are the reason for the bright, vibrant colors of coral reefs. A 1-2°C sustained increase in seawater temperatures is sufficient for bleaching to occur, which turns corals white. If a coral is bleached for a prolonged period of time, death may result. In the Great Barrier Reef, before 1998 there were no such events. The first event happened in 1998 and after that, they began to occur more frequently. Between 2016 and 2020 there were three of them.
Apart from coral bleaching, the reducing pH value in oceans is also a problem for coral reefs because ocean acidification reduces coralline algal biodiversity. The physiology of coralline algal calcification determines how the algae will respond to ocean acidification.
Effects on fisheries
Impacts on marine mammals
Regions and habitats particularly affected
Some effects on marine mammals, especially those in the Arctic, are very direct such as loss of habitat, temperature stress, and exposure to severe weather. Other effects are more indirect, such as changes in host pathogen associations, changes in body condition because of predator–prey interaction, changes in exposure to toxins and emissions, and increased human interactions. Despite the large potential impacts of ocean warming on marine mammals, the global vulnerability of marine mammals to global warming is still poorly understood.
Marine mammals have evolved to live in oceans, but climate change is affecting their natural habitat. Some species may not adapt fast enough, which might lead to their extinction.
It has been generally assumed that the Arctic marine mammals were the most vulnerable in the face of climate change given the substantial observed and projected decline in Arctic sea ice. However, research has shown that the North Pacific Ocean, the Greenland Sea and the Barents Sea host the species that are most vulnerable to global warming. The North Pacific has already been identified as a hotspot for human threats for marine mammals and is now also a hotspot for vulnerability to global warming. Marine mammals in this region will face double jeopardy from both human activities (e.g., marine traffic, pollution and offshore oil and gas development) and global warming, with potential additive or synergetic effects. As a result, these ecosystems face irreversible consequences for marine ecosystem functioning.
Marine organisms usually tend to encounter relatively stable temperatures compared to terrestrial species and thus are likely to be more sensitive to temperature change than terrestrial organisms. Therefore, the ocean warming will lead to the migration of increased species, as endangered species look for a more suitable habitat. If sea temperatures continue to rise, then some fauna may move to cooler water and some range-edge species may disappear from regional waters or experience a reduced global range. Change in the abundance of some species will alter the food resources available to marine mammals, which then results in marine mammals' biogeographic shifts. Furthermore, if a species is unable to successfully migrate to a suitable environment, it will be at risk of extinction if it cannot adapt to rising temperatures of the ocean.
Arctic sea ice decline leads to loss of the sea ice habitat, elevations of water and air temperature, and increased occurrence of severe weather. The loss of sea ice habitat will reduce the abundance of seal prey for marine mammals, particularly polar bears. Sea ice changes may also have indirect effects on animal health due to changes in the transmission of pathogens, impacts on animals' body condition due to shifts in the prey-based food web, and increased exposure to toxicants as a result of increased human habitation in the Arctic habitat.
Sea level rise is also important when assessing the impacts of global warming on marine mammals, since it affects coastal environments that marine mammal species rely on.
Polar bears
Seals
Seals are another marine mammal that are susceptible to climate change. Much like polar bears, some seal species have evolved to rely on sea ice. They use the ice platforms for breeding and raising young seal pups. In 2010 and 2011, sea ice in the Northwest Atlantic was at or near an all-time low and harp seals as well as ringed seals that bred on thin ice saw increased death rates. Antarctic fur seals in South Georgia in the South Atlantic Ocean saw extreme reductions over a 20-year study, during which scientists measured increased sea surface temperature anomalies.
Dolphins
Climate change has had a significant impact on various dolphin species. For example: In the Mediterranean, increased sea surface temperatures, salinity, upwelling intensity, and sea levels have led to a reduction in prey resources, causing a steep decline in the short-beaked common dolphin subpopulation in the Mediterranean, which was classified as endangered in 2003. At the Shark Bay World Heritage Area in Western Australia, the local population of the Indo-Pacific bottlenose dolphin had a significant decline following a marine heatwave in 2011. River dolphins are highly affected by climate change as high evaporation rates, increased water temperatures, decreased precipitation, and increased acidification occur.
North Atlantic right whales
Potential feedback effects
Methane release from methane clathrate
Rising ocean temperatures also have the potential to impact methane clathrate reservoirs located under the ocean floor sediments. These trap large amounts of the greenhouse gas methane, which ocean warming has the potential to release. However, it is currently considered unlikely that gas clathrates (mostly methane) in subsea clathrates will lead to a "detectable departure from the emissions trajectory during this century".
In 2004 the global inventory of ocean methane clathrates was estimated to occupy between one and five million cubic kilometres.
| Physical sciences | Climate change | Earth science |
41501305 | https://en.wikipedia.org/wiki/Telegram%20%28software%29 | Telegram (software) | Telegram Messenger, commonly known as Telegram, is a cloud-based, cross-platform, social media and instant messaging (IM) service. It was originally launched for iOS on 14 August 2013 and Android on 20 October 2013. It allows users to exchange messages, share media and files, and hold private and group voice or video calls as well as public livestreams. It is available for Android, iOS, Windows, macOS, Linux, and web browsers. Telegram offers end-to-end encryption in voice and video calls, and in optional private chats, which Telegram calls Secret Chats.
Telegram also has social networking features, allowing users to post stories, create large public groups with up to 200,000 members, or share one-way updates to unlimited audiences in so-called channels.
Telegram was founded in 2013 by Nikolai and Pavel Durov. Its servers are distributed worldwide with several data centers, while the headquarters is in Dubai, United Arab Emirates. Telegram is the most popular instant messaging application in parts of Europe, Asia, and Africa. It was the most downloaded app worldwide in January 2021, with 1 billion downloads globally as of late August 2021. , registration to Telegram requires either a smartphone or one of a limited number of non-fungible tokens (NFTs) issued in December 2022.
Telegram has more than 950 million monthly active users, with India as the country with the most users.
History
Development
Telegram was launched in 2013 by the brothers Nikolai and Pavel Durov. Previously, the pair founded the Russian social network VK, which they left in 2014, saying it had been taken over by the government. Pavel sold his remaining stake in VK and left Russia after resisting government pressure. Nikolai created the MTProto protocol that is the basis for the messenger, while Pavel provided financial support and infrastructure through his Digital Fortress fund. Telegram Messenger states that its end goal is not to bring profit, but it is not structured as a non-profit organization.
Telegram is registered as a company in the British Virgin Islands and as an LLC in Dubai. It does not disclose where it rents offices or which legal entities it uses to rent them, citing the need to "shelter the team from unnecessary influence" and protect users from governmental data requests. After Pavel left Russia in 2014, he was said to be moving from country to country with a small group of computer programmers consisting of 15 core members.
While a former employee of VK said that Telegram had employees in Saint Petersburg, Pavel said the Telegram team made Berlin, Germany, its headquarters in 2014, but failed to obtain German residence permits for everyone on the team and moved to other jurisdictions in early 2015. Since 2017, the company has been based in Dubai. Its data centers are spread across a complex corporate structure of shell companies in various jurisdictions to avoid compliance with government subpoenas. The company says this is done "to protect the data that is not covered by end-to-end encryption". Telegram's FAQ page says it does not process any requests related to illegal content in chats and group chats, and that "to this day, we have disclosed 0 bytes of user messages to third parties, including governments". However, according to Pavel, Telegram has disclosed data for 203 legal requests from the Brazilian government from Q1 to Q3 of 2024, and a total of 6,992 legal requests from India, its largest market, during the same period. Users can check how many legal requests from their country have been processed by Telegram using the official transparency bot.
Usage
In October 2013, Telegram announced that it had 100,000 daily active users.
On 24 March 2014, Telegram announced that it had reached 35 million monthly users and 15 million daily active users. In October 2014, South Korean government surveillance plans drove many of its citizens to switch to Telegram from the Korean app KakaoTalk. In December 2014, Telegram announced that it had 50 million active users, generating 1 billion daily messages, and that it had 1 million new users signing up on its service every week, traffic doubled in five months with 2 billion daily messages. In September 2015, Telegram announced that the app had 60 million active users and delivered 12 billion daily messages.
In February 2016, Telegram announced that it had 100 million monthly active users, with 350,000 new users signing up every day, delivering 15 billion messages daily. In December 2017, Telegram reached 180 million monthly active users. By March 2018, that number had doubled, with Telegram reaching 200 million monthly active users.
On 14 March 2019, Pavel said that "3 million new users signed up for Telegram within the last 24 hours." He did not specify what prompted this flood of new sign-ups, but the period matched a prolonged technical outage experienced by Facebook and its family of apps, including Instagram. According to the US Securities and Exchange Commission, as of October 2019, Telegram had 300 million monthly active users worldwide.
On 24 April 2020, Telegram announced that it had reached 400 million monthly active users.
On 8 January 2021, Pavel announced in a blog post that Telegram had reached "about 500 million" monthly active users. In August, TechCrunch reported that India was Telegram's largest market, with a 22% share of total installs coming from the region. Telegram then gained over 70 million new users as a result of an outage which affected Facebook and its affiliates on 5 October 2021.
On 19 June 2022, Telegram announced that it had reached 700 million monthly active users.
In July 2023, Telegram surpassed 800 million monthly active users, later reaching 900 million in March 2024 and 950 million in July 2024.
Features
Messaging
To start using Telegram, a user must sign-up with their phone number or an anonymous +888 number purchased from the Fragment blockchain platform. Changing the phone number in the app will automatically reassign the user's account to that number without the need to export data or notify their contacts. Phone numbers are hidden by default with only a user's contacts being able to see them. Sign-ups can only be done via an Android or iOS device.
Upon signing up, messages sent and received by the user are tied to their number and a custom username, not the device. Any Telegram content is synced between the user's logged-in devices automatically through cloud storage, except for device-specific secret chats. By default, any account that is inactive for 6 months is automatically deleted, though the period can be shortened or extended up to 18 months through the Settings menu. Telegram allows groups, bots and channels with a verified social media or Wikipedia page to be verified, but not individual user accounts.
Messages can contain formatted text, media, files up to 2 GB (4 GB with Premium), locations and audio or video messages recorded in-app. Telegram messages in private chats can be edited up to 48 hours after they were sent with an "edited" icon appearing to reflect changes, as well as deleted for both sides without a trace. Users have the option to delete messages and whole chats for both themself and other participants. Chats can be exported to preserve them via Telegram's Desktop client, although the saved data cannot be imported back into the user's account.
Users can however import chat history, including both messages and media, from WhatsApp, Line, and KakaoTalk due to data portability, either making a new chat to hold the messages or adding them to an existing one.
As users can be logged into many devices at once, starting to type a message on one of them will create a "cloud draft" that syncs with others, so typing can be started on a phone and finished on a laptop, for example.
Any message in any chat can be translated by opening the context menu. Premium users have the option to translate the whole chat with one click. Users can hide the translate button for messages written in specific languages.
Reactions can be used to respond to a message with emoji, with Premium users having access to more reaction choices and the ability to leave more reactions per message. Reactions are always on in private chats and can be enabled by admins in groups and channels with the ability to allow or exclude specific reactions. Reaction emoji play an animation with special effects when sent.
Users can also send stickers, which can be static, animated or video. Sticker packs are made by Telegram designers as well as regular users and can be shared via links. They use the WebP or WebM format and do not require special software to create or upload. Some stickers feature full-screen effects that play out when first sent or when tapped.
Users can schedule messages to send at a particular time or when their conversation partner comes online, as well as choose to send a message "without sound" without a notification. Messages from private chats can be forwarded, with an option to hide the original sender's identity or to hide captions from media messages. Forwarded messages also maintain reply formatting, able to show which messages in a thread are replying to others. Any user can also send a message to a special "Saved Messages" chat as a form of bookmarking them. The contents of the chat are only visible to the user.
Chats can be sorted into folders to organize them with preset options like "Unread" and "Muted? or custom separations such as "Work" and "Family". Premium users have the ability to set any chat folder as the default screen in the app while regular users will always see the full chat list when first opening the app.
Users have the option to start a one-on-one, end-to-end-encrypted "Secret Chat", which remains accessible only on the device where it was started and self-destructs upon logging out. Secret Chats restrict screenshotting from Android devices and warn when one is taken from an iOS device, while also hiding the chat contents from the final image. Secret Chats support perfect forward secrecy and switch encryption keys after a key has been used 100 times or a week has passed. Secret Chats are only available on Android, iOS and macOS clients.
Both in Secret and regular chats, messages can self-destruct after they are read, disappearing for all parties after a period set by the user, ranging from 1 day to 1 year.
Groups and channels
Telegram users can create and join groups and channels. Groups are large multi-user chats that support up to 200,000 members and can be public or private. Users can freely join public chats and find them using the in-app search function, while private chats require an invitation. They support flexible admin rights and can use bots for moderation to prevent spam and unwanted activity. Groups can be split into topics, effectively creating subgroups dedicated to various subjects with separate settings for each.
Admins can choose to hide the list of members in a group, as well as post anonymously themselves. Similarly, groups and channels can have content protection enabled, which prevents screenshots, forwarding and downloading of media. Ownership of channels and groups can be transferred to one of the admins if the owner wishes to give up their rights.
Groups support threaded replies, where bringing up the context menu on a message allows one to open a screen with a thread of replies made to that message and the subsequent ones in the thread. Specific users can be tagged in the group by adding @username to a message, where "username" is that particular user's username.
Groups and channels also support polls, which can be open or anonymous and can support multiple choices. When forwarded, polls retain the answer data and any votes cast in other chats will count toward the overall total.
Channels are one-way feeds where the channel owner or admins can post content while followers can only read, react and comment, if comments have been enabled. Channels can be created for broadcasting messages to an unlimited number of subscribers. The list of those who subscribe to a channel can only be seen by its admins. Posting in the channel is anonymous, though admins can choose to add signatures to their posts. Channels offer detailed statistics on view counts, user growth and interactions, also visible only to admins.
Channel owners are able to use Telegram to create giveaways, randomly awarding channel members with prizes such as Telegram Premium subscriptions to their followers, based on certain criteria. Users with a Telegram Premium subscription have a number of "boosts" that they can give to channels, which allow the channel to "level up" and unlock features, such as the ability to customize messages or post stories as the channel.
In December 2019, Bloomberg News moved their messenger-based newsletter service from WhatsApp to Telegram after the former banned bulk and automated messaging. Other news services with official channels on the platform include the Financial Times, Business Insider and The New York Times.
Channels have also been used by governments and heads of state. Notable examples include Volodymyr Zelenskyy and Emmanuel Macron. Channels have been used by journalists in oppressive regimes to establish independent news networks.
Games
Telegram also provides an open API for the creation of custom bots which can perform various tasks, integrate other services into Telegram chats, or work as mini apps or games. Most of them work on the 8XR game engine.
Video and voice calls
Since 2017, Telegram users have been able to initiate one-on-one calls in private chats. Calls are end-to-end encrypted and prioritize peer-to-peer connections. Video calls were introduced in August 2020. According to Telegram, there is a neural network working to learn various technical parameters about a call to provide better quality of service for future uses.
Telegram added group voice chats in December 2020 and group video chats in June 2021. Group voice and video chats support picture-in-picture video, as well as sharing one's screen, creating a recording of the call, noise suppression and selective muting. In channels, users can start a livestream, that is able to integrate with third-party apps such as OBS Studio and XSplit.
Once launched, a group voice chat will remain active and open to all group members until an admin specifically closes it.
Privacy and security features
By default, logging into Telegram requires either an SMS message sent to the registered number or a code message sent to one of the active sessions on another device. Users have the option to set a two-step verification password and add a recovery email. In late 2022, options to Sign in with Apple and Sign In with Google or with an email address were added. Whenever a new device successfully logs in to a user's account, a special service notification is sent and a login alert is displayed in the chat list of their other devices.
In the Privacy and Security submenu of Settings, users have the option to hide their "Last Seen" status, which reflects the last time the user opened a Telegram app. Hiding the status obfuscates the exact time of the user being online and hides the statuses of other people respectively. Similarly, users can hide their phone number and profile photo from people based on categories such as Non-Contacts or by adding exceptions. When a user chooses to hide their profile photo, they have an option of setting an alternative "Public Profile Picture" that will be shown instead.
In the same menu, users can restrict the circle of people who can call them or invite them to groups and channels, while Premium users also have the option to restrict who can send them voice messages.
The Devices submenu shows all of the active devices on a user's account and allows them to remotely log out from those devices.
Data and storage settings
Telegram clients have the ability to turn off media autoplay and automatic downloads for both WiFi and mobile data, adjusting them for media type and size. Auto download settings can also be applied based on chat type such as group, channel or private.
Cache settings can be changed to automatically clear the cache once it reaches a certain size or a certain time passes. The interface shows users a visual representation of their storage usage and also lets them sort their cached media by size to clear specific items.
Bots
In June 2015, Telegram launched a platform for third-party developers to create bots. Bots are Telegram accounts operated by programs. They can respond to messages or mentions directly or can be invited into groups, and are able to perform tasks, integrate with other programs and host mini apps. Bots can accept online payments made with credit cards or Apple Pay. The Dutch website Tweakers reported that an invited bot can potentially read all group messages, when the bot controller changes the access settings silently at a later point in time. Telegram pointed out that it considered implementing a feature that would announce such a status change within the relevant group.
There are also inline bots, which can be used from any chat screen. To activate an inline bot, a user must type the bot's username and a query in the message field. The bot then will offer its content. The user can choose from that content and send it within a chat. Certain approved bots are also able to integrate into the attachment menu, making them accessible in any chat.
Bots can handle transactions provided by Paymentwall, Yandex.Money, Stripe, Ravepay, Razorpay, QiWi and Google Pay for different countries. Bots power Telegram's gaming platform, which utilizes HTML5, so games are loaded on-demand as needed, like ordinary webpages. Games work on iPhone 4 and newer, and on Android 4.4 devices and newer.
People can use Internet Of Things (IoT) services with two-way interaction via IFTTT implemented within Telegram.
In April 2021, the Payments 2.0 upgrade enabled bot payments within any chat, using third-party services such as Sberbank, Tranzoo, Payme, CLICK, LiqPay and ECOMMPAY to process the credit card information.
In February 2018, Telegram launched their social login feature to its users, named Telegram Login. It features a website widget that could be embedded into websites, allowing users to sign into a third party website with their Telegram account. The gateway sends the user's Telegram name, username, and profile picture to the website owner, while the user's phone number remains hidden. The gateway is integrated with a bot, which is linked with the developer's specific website domain.
In June 2021, an update introduced a new bot menu where users can browse and send commands while in a chat with a bot.
In April 2022, bots gained support for customized interfaces and inline page loading. Interfaces can be adjusted to match the app's theme even if it changes while interacting with the bot.
In October 2024, Telegram added increased messaging limits for bots, allowing bots to send up to 1000 messages per second to their users. Messages beyond the free limit of 30 per second are paid for using Telegram Stars.
Telegram introduced affiliate programs in December 2024, which allow developers to create an affiliate program for their bot or mini app. Any Telegram user can join the affiliate program, and are rewarded for referring others to the bot or mini app by receiving a commission from purchases made by the people they referred.
Stickers, emoji, reactions and effects
Telegram has more than 40,000 stickers. Stickers are cloud-based, high-resolution images intended to provide more expressive emoji. When typing in an emoji, the user is offered to send the respective sticker instead. Stickers come in collections called "packs", and multiple stickers can be offered for one emoji. Telegram comes with one default sticker pack, and users can install additional sticker packs provided by third-party contributors.
Sticker sets installed from one client become automatically available to all other clients. Sticker images use WebP file format, which is better optimized to be transmitted over the internet. The Telegram clients also support animated emoji. In January 2022, video stickers were added, which use the WebM file format and do not feature any software requirements to create.
In August 2019, Telegram introduced animated emoji, larger versions of familiar emoji with unique animations. In September 2021, Telegram added interactive emoji, a type of animated emoji which also play a fullscreen effect in the chat. These kinds of effects were later used for Premium Stickers in June 2022 and for message effects in May 2024.
In August 2022, Telegram launched an emoji platform where users could upload their own custom emoji, either in animated or static versions. While any user can upload custom emoji to the platform, the use of custom emoji in chats is only available to users with Telegram Premium.
Reactions were first added to Telegram in 2021 and expanded to include more emoji options for Premium users. In September 2022, Telegram gave free users access to dozens of reactions, even some that were only previously available to Premium subscribers. In order to accommodate the new reactions, the reaction panel was expanded and redesigned.
People Nearby and Groups Nearby
People Nearby can help users meet new friends by turning on phone GPS location and opting-in contacts and through Groups Nearby people can create a local group by adding location data to groups.
It is planned that People Nearby will be sunsetted in the future.
Stories
Similar to other social platforms, Telegram users can post stories, a type of short-form content. Telegram stories have several distinctive features, like a dual-camera mode, extra privacy settings, the ability to edit stories after posting them, as well as to rewind and fast-forward them while watching.
Premium features
Telegram Premium was launched on 19 June 2022 with regional pricing. The optional paid subscription gives users increased limits in the app, such as larger file uploads, faster download speeds, unlimited voice message transcription, as well as numerous other increases such as the number of pinned chats and folders. Premium users have access to extra stickers, emoji, reactions, and customization features like a special badge and the ability to change the look of their messages in chats. Premium users get access to additional settings, like instant chat translation, and the ability to restrict who can send them voice messages.
As of 2023, Telegram Premium can be acquired via in-app purchases facilitated by Apple and Google, directly via Telegram's @PremiumBot, or with cryptocurrency on the Fragment platform. Users are able to purchase a subscription for themselves, or purchase a subscription for someone else to send as a gift. Premium subscriptions can also be won through official Channel Giveaways, in which Telegram channels pre-purchase a specific number of Premium subscriptions that are randomly given away to their subscribers.
Related platforms
People can use their Telegram accounts to author articles on Telegraph – a minimalistic text editor and publisher. While articles on Telegraph can be published anonymously, tying them to one's account allows one to check their view count and edit them later. Telegraph natively supports Instant View, a feature which lets users read full articles in the chat with no load time and without opening an external browser.
When an article is first published, the URL is generated automatically from its title. Non-Latin characters are transliterated, spaces are replaced with hyphens, and the date of publication is added to the address. For example, an article titled "Telegraph (blog platform)" published on 17 November would receive the URL /Telegraph-blog-platform-11-17.
Text formatting options are also minimal: two levels of headings, single-level lists, bold, italics, quotes, and hyperlinks are supported. Authors could upload images and videos to the page, with a limit of 5 MB, however, it has been disabled since September 2024. When an author adds links to YouTube, Vimeo, or Twitter, the service allows you to embed their content directly in the article.
In February 2018, Telegram launched their social login feature to its users, named Telegram Login. It features a website widget that could be embedded into websites, allowing users to sign into a third party website with their Telegram account. The gateway sends users' Telegram name, username, and profile picture to the website owner, while users' phone number remains hidden. The gateway is integrated with a bot, which is linked with the developer's specific website domain.
In July 2018, Telegram introduced their online authorization and identity-management system, Telegram Passport, for platforms that require real-life identification. It asks users to upload their own official documents such as passport, identity card, driver license, etc. When an online service requires such identification documents and verification, it forwards the information to the platform with the user's permission. Telegram stated that it does not have access to the data, while the platform will share the information only with the authorized recipient. However, the service was criticised for being vulnerable to online brute-force attacks.
In December 2020, Telegram launched a Bugs and Suggestions platform, where users can submit bug reports and suggestion cards for new features. Others can then vote and comment on the cards.
In October 2024, Telegram launched a verification platform, called Telegram Gateway, allowing third-party services to authenticate their users by sending verification codes via Telegram.
Architecture
Privacy
For encrypted chats (branded as Secret Chats), Telegram uses a custom-built symmetric encryption scheme called MTProto. The protocol was developed by Nikolai Durov and other developers at Telegram and, as of version 2.0, is based on 256-bit symmetric AES encryption, 2048-bit RSA encryption and Diffie–Hellman key exchange.
MTProto 1.0 was deprecated in favor of MTProto 2.0 in December 2017, which was deployed in Telegram clients as of v4.6.
Version 2.0 was proven formally correct in December 2020 by a team from the University of Udine, Italy. The team reviewed the protocol after realizing that they could only find in-depth verifications done of version 1.0, where most criticisms were levied. They used ProVerif, a verifier based on the symbolic Dolev-Yao model. In the published paper, they "provide a fully automated proof of the soundness of MTProto 2.0’s protocols for authentication, normal chat, end-to-end encrypted chat, and re-keying mechanisms with respect to several security properties, including authentication, integrity, confidentiality and perfect forward secrecy...MTProto 2.0 is assumed to be a perfect authenticated encryption scheme (IND-CCA and INT-CTXT)."
However, the team also stated that because all communication, including plaintext and ciphertext, passes through Telegram servers, and because the server is responsible for choosing Diffie–Hellman parameters, the "server should not be considered as trusted." They also concluded that a man-in-the-middle attack is possible if users fail to check the fingerprints of their shared keys. Finally, they qualified their conclusion with the caveat that "properties need to be formally proved in order to deem MTProto 2.0 definitely secure. This proof cannot be done in a symbolic model like ProVerif's, but it can be achieved in a computational model, using tools like CryptoVerif or EasyCrypt."
Servers
As with most instant messaging protocols, Telegram uses centralized servers. Telegram Messenger LLP has servers in a number of countries throughout the world to improve the response time of their service. Telegram's server-side software is closed-source and proprietary. Pavel Durov said that it would require a major architectural redesign of the server-side software to connect independent servers to the Telegram cloud.
For users who signed in from the European Economic Area (EEA) or United Kingdom, the General Data Protection Regulations (GDPR) are supported by storing data only on servers in the Netherlands, and designating a London-based company as their responsible data controller.
Clients
Telegram has various client apps, some developed by Telegram Messenger LLP and some by the community. Most of them are free and open-source and released under the GNU General Public Licence or 3. The official clients support sending any file format extensions. The built-in media viewer supports common media formats – JPEG, PNG, WebP for images and H.264 and HEVC in videos in MP4 container and MP3, FLAC, Vorbis, Opus and AAC for audio.
In 2021, the Telegram team announced a direct build of its Android app. Telegram for Android is available directly from the Telegram website. It is automatically updated and will most likely get new versions faster than the apps in the Play Store and App Store. A distinctive feature of this version is the ability to view channels/groups on a specific topic without censorship, which cannot be viewed from an app distributed from Google Play or the Apple Store due to their policies.
Common specifications:
No cloud backup option for secret chat
APIs
Telegram has public APIs with which developers can access the same functionality as Telegram's official apps to build their own instant messaging applications. In February 2015, creators of the unofficial WhatsApp+ client released the Telegram Plus app, later renamed to Plus Messenger, after their original project got a cease-and-desist order from WhatsApp. In September 2015, Samsung released a messaging application based on these APIs.
Telegram also offers an API that allows developers to create bots, which are accounts controlled by programs. Such bots are used, among other things, to emulate and play old games in the app and inform users about vaccine availability for COVID-19.
In addition, Telegram offers functions for making payments directly within the platform, alongside an external service such as Stripe.
Business
The company was initially supported by founding CEO Pavel Durov's personal funds after the sale of his stake in VK. In January 2018, it launched a private placement and collected $1.7 billion from investors such as Kleiner Perkins, Sequoia Capital, and Benchmark. After the shutdown of the TON project, the company needed to repay the investors the money that was not spent on its development during 2018 and the beginning of 2019, while the project was active.
On 15 March 2021, Telegram conducted a five-year public bonds placement worth $1 billion. The funding was required to cover the debts amounting to $625.7 million, including $433 million to investors who bought futures for Gram tokens in 2018 and included purchasers such as David Yakobashvili. On 23 March, Telegram sold additional bonds worth $150 million to the Abu Dhabi Mubadala Investment Company and Abu Dhabi Catalyst Partners.
A day later, the Mubadala Investment Company stated that Russia's sovereign wealth fund participated in its deal undisclosed through the Russia-UAE joint investment platform to buy convertible bonds. A Telegram spokesperson stated: "RDIF is not in the list of investors we sold bonds to. We wouldn't be open to any transaction with this fund" and "[t]he funds that did invest, including Mubadala, confirmed to us that RDIF was not among their LPs [limited partners]." According to the contract, the holders of the bonds will be provided with an option to convert them to shares at a 10% discount if the company conducts an open IPO.
Durov stated that the move aimed to "enable Telegram to continue growing globally while sticking to its values and remaining independent". According to press reports, prior to the bonds placement, Durov had rejected an investment offer for a 5–10% stake in the company as well as several undisclosed ones, valuing the company in a $30–40 billion range. In March 2024, Telegram sold an additional $330 million in bonds. Durov said the bond sale "will further solidify our position as an independent platform that is able to challenge the 'Goliaths' of our industry".
Advertising and monetization
Telegram has stated that the company will never serve advertisements in private chats. In late 2020, Durov announced that the company was working on its own ad platform, and would integrate non-targeted ads in public one-to-many channels, that already were selling and displaying ads in the form of regular messages. Ads from Telegram's "Sponsored Messages" platform began to appear in channels with more than 1000 followers in October 2021.
In late 2020, Durov announced that Telegram will consider adding paid features aimed at enterprise clients. According to him, these features will require more bandwidth and the added cost will be covered by the feature prices, in addition to covering some of the costs incurred by regular users.
In June 2024, Telegram launched Telegram Stars to facilitate in-app purchases of digital goods and services, in compliance with policies from the App Store and Play Store. Following the launch of Stars, Telegram released several updates to their functionality, such as allowing Stars to be used to unlock media in channels or to buy gifts for other users.
In December 2024, Telegram CEO Pavel Durov announced that Telegram had reached profitability, due to significant growth in Premium subscriptions and Telegram ad sales, assisted by the other monetization features launched throughout 2024.
TON Telegram Open Network
In 2017, in an attempt to monetize Telegram without advertising, the company began the development of a blockchain platform dubbed either "The Open Network" or "Telegram Open Network" (TON) and its native cryptocurrency "Gram". The project was announced in mid-December 2017 and its 132-page technical paper became available in January 2018. The codebase behind TON was developed by Pavel Durov's brother Nikolai Durov, the core developer of Telegram's MTProto protocol. In January 2018, a 23-page white paper and a detailed 132-page technical paper for TON blockchain became available.
Durov planned to power TON with the existing Telegram user base, and turn it into the largest blockchain and a platform for apps and services akin to a decentralized WeChat, Google Play, and App Store. Besides, the TON had the potential to become a decentralized alternative to Visa and MasterCard due to its ability to scale and support millions of transactions per second. In January and February 2018, the company ran a private sale of futures contracts for Grams, raising around $1.7 billion. No public offering took place.
The development of TON took place in a completely isolated manner, and the release was postponed several times. The test network was launched in January 2019. The launch of the TON main network was scheduled for 31 October. On 30 October, the U.S. Securities and Exchange Commission obtained a temporary restrictive order to prevent the distribution of Grams to initial purchasers; the regulator considered the legal scheme employed by Telegram as an unregistered securities offering with initial buyers acting as underwriters.
The judge hearing the Telegram v. SEC case, P. Kevin Castel, ultimately agreed with the SEC's argument and kept the restrictions on Gram distribution in force. The ban applied to non-U.S.-based purchasers as well, because Telegram could not prevent the re-sale of Grams to U.S. citizens on a secondary market, as the anonymity of users was one of the key features of TON. Following that, Durov announced the end of Telegram's active involvement with TON. On 26 June, the judge approved the settlement between Telegram and SEC. The company agreed to pay an $18.5 million penalty and return $1.22 billion to Gram purchasers. In March 2021, Telegram launched a bonds offering to cover the debt and fund further growth of the app.
Criticism
Due to Telegram's mixed nature as both a private communication method and a social media-like platform with mass groups and channels, along with its minimal restrictions on content with only calls to violence, illegal forms of pornography and scamming forbidden, the platform has been used by organizations and large groups for recruitment and spreading of their agenda. Organized use of the app has been linked to pro-democracy protests in Belarus, Russia, Hong Kong and Iran, as well as to the dissemination of state propaganda and violent rhetoric in oppressive regimes, the promotion of extremist views, and the digitalization of services provided by government entities and private businesses.
The app has been criticized by numerous research institutions and internet monitoring bodies due to violent organizations like ISIS, Proud Boys and the Myanmar junta using the app to communicate, both privately between members and publicly through channel posts. While Telegram made substantial efforts to ban illegal content such as child abuse and pro-terrorist channels, including a partnership with Europol to eliminate IS presence on the platform, communities of far-left and far-right, antivaxx, Antifa and extremist users are still found on the app. Such content is usually linked to Telegram allowing misinformation on the platform as, according to founder and CEO Pavel Durov, "conspiracy theories only strengthen each time their content is removed by moderators".
In September 2024, Telegram announced that it will begin to hand over users' IP addresses and phone numbers to authorities who have search warrants or other valid legal requests.
Use by militant groups
In September 2015, in response to a question about the use of Telegram by Islamic State (ISIS), Pavel Durov stated: "I think that privacy, ultimately, and our right for privacy is more important than our fear of bad things happening, like terrorism." Durov sarcastically suggested to ban words because terrorists use them for communication.
ISIS has used Telegram for recruiting attempts with some cells recommending the app to their followers. In France, initial investigations of a terrorist act revealed the perpetrators used Telegram to communicate, though follow-up research suggested that the extent of the app's use was "unclear". Beginning in December 2016, Telegram began publishing daily moderation statistics regarding terror-related content in an official channel named @ISISwatch. Following efforts by Telegram to remove ISIS-related content from the platform, the terrorist organization reportedly moved its recruitment groups to the dark web, with US officials citing the app's purging of terrorist content as particularly effective at deplatforming ISIS.
In 2023, Saudi Arabia's Global Center for Combating Extremist Ideology (Etidal) reported that in collaboration with Telegram, over 59 million pieces of extremist content had been removed from the platform since 2022.
Throughout 2020, the Islamic Revolutionary Guard Corps used groups and channels to dox Iraqi and Iranian citizens, while sharing propagandistic posts on the platform. After the 2021 coup d'état in Myanmar, Telegram channels were used by the junta to spread propaganda and organize misinformation campaigns against pro-democracy groups. In response to the Office of the United Nations High Commissioner for Human Rights call to prevent such abuse, Telegram reportedly banned 13 accounts related to or supporting the Myanmar military.
On 26 April 2023, Telegram was temporarily suspended throughout Brazil, and the company was fined per day for not complying with a Federal Police investigation into neo-Nazi activities on the platform. The company only partially fulfilled a court request for personal data on two anti-Semitic Telegram groups, which authorities considered an intentional lack of cooperation. The decision was made after a series of violent school attacks, with at least one incident being linked to exchanges on an anti-Semitic group. The Telegram CEO then said that the requested data was technologically impossible to obtain. A federal court lifted the suspension three days later, but upheld the daily fine. Twelve days later, Telegram told its users that the Brazilian Congressional Bill No. 2630 against online disinformation, which was about to be approved, would end freedom of speech in the country.
Far-right and white supremacist communities on Telegram have spread videos of the Christchurch and Halle shootings in groups and channels after the original livestreams were taken down by Twitch. British far-right publication TR.news, following multiple deplatformings, launched a Telegram channel to spread its posts and the Institute for Strategic Dialogue reported that Irish far-right groups grew substantially between 2019 and 2020. However, research from the Oxford University suggests that, due to Telegram not using sorting algorithms in its search function, many such groups remain obscure and small while select others receive a lot of attention.
Illegal pornography
Telegram has been used to distribute illegal pornography, including child pornography. Telegram's internal reporting system has an option to report content that contains child abuse, including specific messages in groups and channels. The company has a verified channel called "Stop Child Abuse", where daily statistics on the number of groups and channels banned for sharing illegal materials are posted. It also provides an email address dedicated to reports of content related to child abuse.
In January 2021, Macedonian media outlets reported that a now-banned Telegram group, with more than 7,000 members, titled "Public Room" ("Јавна соба") was used to share nude photos of women, often young teenage girls. Along with the shared photographs, anonymous accounts shared private information of the women, including phone numbers and social media profiles, encouraging members of the group to contact the women and ask for sexual favors.
This was done without prior agreement or knowledge of the women, causing intense public backlash and demand for the group to be shut down. The President of North Macedonia Stevo Pendarovski, along with the Prime Minister of North Macedonia Zoran Zaev, demanded an immediate reaction from Telegram and threatened to completely restrict access to the app in the country if no actions were taken. The group was banned after reports from users and media, although no public statement was made.
An investigation published in August 2024 by the BBC showed that Telegram has not responded to requests to join the US–based National Center for Missing & Exploited Children (NCMEC), nor from the UK–based Internet Watch Foundation, both non–profit NGOs. Telegram responded to the BBC's reporting, stating that Telegram "proactively moderates harmful content on its platform including child abuse material", and that its moderation is "within industry standards and constantly improving".
In August 2024, Journalist Ko Narin of Hankyoreh exposed Telegram chats of teenagers who used AI to deepfake images of their classmates and teachers for porn.
26 November 2024 BBC and Radiofarda discovered Iranian regime hosting CP posting channels.
Bot abuse
Volodymyr Flents, the chairman of the public organization "Electronic Democracy", announced on 11 May 2020 that a Telegram bot appeared on the Web, which sold the personal data of Ukrainian citizens. It is estimated that the bot contains data from 26 million Ukrainians registered in the Diia application. However, deputy prime minister and minister of digital transformation Mykhailo Fedorov denied any data from the app being leaked. The criminal activity of 25 people was confirmed and copies of 30 databases were seized.
In 2020, a Telegram bot was blocked by Apple after posting deepfake pornography. In the same year, Telegram reportedly banned more than 350,000 bots and channels, including those that contained child abuse and terrorism-related content.
In 2021, a bot was found selling leaked phone numbers from Facebook.
Fraudulent jobs
Telegram has received criticism for its failure to curb fraud on the platform. The most common mode of fraud involves scammers sending messages to unsuspecting users, offering part-time online jobs which comprise a series of tasks. Scammers employ a variety of confidence tricks to entice users into completing "prepaid tasks" in which users deposit money into scammers' accounts with the expectation of receiving high returns.
In July 2023, Hyderabad Police uncovered a fraud wherein 15,000 Indian citizens were duped out of in less than a year, all related to "prepaid tasks" on Telegram. A cybercrime police investigation of the money trail revealed that the fraud originated from China and the money was laundered by mules through cryptocurrency wallets.
In September 2023, the Singapore Police Force stated that more than 6,600 Singaporeans had lost over to prepaid job scams on Telegram and WhatsApp since the start of the year.
Copyright infringement
In March 2024, a judge of Spain's Audiencia Nacional ordered the temporary blocking of Telegram in Spain. The order came following a complaint from media organizations —Mediaset, Atresmedia and Movistar Plus+— saying that the app allowed users to share copyrighted content without their consent. A few days later, following repeated criticism, the same judge suspended his order until the police issue a report on the consequences that this measure would have for users. Finally, the judge annulled the order, considering it "disproportionate".
Drug trade
In recent years, Telegram has become more popular for the purpose of buying and selling illicit drugs. In 2024, Sociology Compass released a paper exploring this trend in drug trade.
User numbers
In August 2024, an EU probe was launched into Telegram to determine whether the platform breached EU digital rules by failing to provide accurate user numbers. Telegram said in February 2024 that it had 41mn users in the EU. Under the EU's Digital Services Act (DSA), Telegram was supposed to provide an updated number in August but failed to do so, only declaring it had "significantly fewer than 45mn average monthly active recipients in the EU".
Reception
Channels have been used by celebrities such as Arnold Schwarzenegger and politicians: President of France Emmanuel Macron, former President of Brazil Jair Bolsonaro, Turkish President Recep Tayyip Erdoğan, President of Moldova Maia Sandu, President of Ukraine Volodymyr Zelenskyy, former Mexican President Andrés Manuel López Obrador, former Singaporean Prime Minister Lee Hsien Loong, US President Donald Trump, President of Uzbekistan Shavkat Mirziyoyev, former Taiwan President Tsai Ing-wen, Ethiopian Prime Minister Abiy Ahmed, Israeli Prime Minister Benjamin Netanyahu and others.
Security
Telegram's security model has received notable criticism by cryptography experts. They criticized how, unless modified first, the default general security model stores all contacts, messages and media together with their decryption keys on its servers continuously; and that it does not enable end-to-end encryption for messages by default. Pavel Durov has argued that this is because it helps to avoid third-party unsecured backups, and to allow users to access messages and files from any device. Criticisms were also aimed at Telegram's use of a custom-designed encryption protocol.
In December 2020, a study titled "Automated Symbolic Verification of Telegram's MTProto 2.0" was published, confirming the security of the updated MTProto 2.0 and reviewing it while pointing out several theoretical vulnerabilities. The paper provides "fully automated proof of the soundness of MTProto 2.0's authentication, normal chat, end-to-end encrypted chat, and re-keying mechanisms with respect to several security properties, including authentication, integrity, confidentiality and perfect forward secrecy" and "proves the formal correctness of MTProto 2.0". This partially addresses the concern about the lack of scrutiny while confirming the formal security of the protocol's latest version.
The desktop clients, excluding the macOS client, do not feature options for end-to-end encrypted messages. When the user assigns a local password in the desktop application, data is also locally encrypted. Telegram has defended the lack of ubiquitous end-to-end encryption by saying that online-backups that do not use client-side encryption are "the most secure solution currently possible".
In May 2016, critics disputed claims by Telegram that it is "more secure than mass market messengers like WhatsApp and Line", as WhatsApp claims to apply end-to-end encryption to all of its traffic by default and uses the Signal Protocol, which has been "reviewed and endorsed by leading security experts", while Telegram does neither and stores all messages, media and contacts in their cloud. Since July 2016, Line has also applied end-to-end encryption to all of its messages by default, though it has also been criticized for being susceptible to replay attacks and the lack of forward secrecy between clients.
In 2013, an author on the Russian programming website Habr discovered a weakness in the first version of MTProto that would allow an attacker to mount a man-in-the-middle attack and prevent the victim from being alerted by a changed key fingerprint. The bug was fixed on the day of the publication with a $100,000 payout to the author and a statement on Telegram's official blog.
On 26 February 2014, the German consumer organization Stiftung Warentest evaluated several data-protection aspects of Telegram, along with other popular instant-messaging clients. Among the aspects considered were: the security of the data transmission, the service's terms of use, the accessibility of the source code, and the distribution of the app. Telegram was rated 'problematic' () overall. The organization was favorable to Telegram's secure chats and partially free code but criticized the mandatory transfer of contact data to Telegram's servers and the lack of an imprint or address on the service's website. It noted that while the message data is encrypted on the device, it could not analyze the transmission due to a lack of source code.
The Electronic Frontier Foundation (EFF) listed Telegram on its "Secure Messaging Scorecard" in February 2015. Telegram's default chat function received a score of 4 out of 7 points on the scorecard. It received points for having communications encrypted in transit, having its code open to independent review, having the security design properly documented, and having completed a recent independent security audit. Telegram's default chat function missed points because the communications were not encrypted with keys the provider did not have access to, users could not verify contacts' identities, and past messages were not secure if the encryption keys were stolen. Telegram's optional secret chat function, which provides end-to-end encryption, received a score of 7 out of 7 points on the scorecard. The EFF said that the results "should not be read as endorsements of individual tools or guarantees of their security", and that they were merely indications that the projects were "on the right track".
In December 2015, two researchers from Aarhus University published a report in which they demonstrated that MTProto 1.0 did not achieve indistinguishability under chosen-ciphertext attack (IND-CCA) or authenticated encryption. The researchers stressed that the attack was of a theoretical nature and they "did not see any way of turning the attack into a full plaintext-recovery attack". Nevertheless, they said they saw "no reason why [Telegram] should use a less secure encryption scheme when more secure (and at least as efficient) solutions exist". The Telegram team responded that the flaw does not affect message security and that "a future patch would address the concern". Telegram 4.6, released in December 2017, supports MTProto 2.0, which now satisfied the conditions for IND-CCA. MTProto 2.0 is seen by qualified cryptographers as a vast improvement to Telegram's security.
In April 2016, the accounts of several Russian opposition members were hijacked by intercepting the SMS messages used for login authorization. In response, Telegram recommended using the optional two-factor authentication feature. In May 2016, the Committee to Protect Journalists and Nate Cardozo, senior staff attorney at Electronic Frontier Foundation, recommended against using Telegram because of "its lack of end-to-end encryption [by default] and its use of non-standard MTProto encryption protocol, which has been publicly criticized by cryptography researchers, including Matthew Green".
On 2 August 2016, Reuters reported that Iranian hackers compromised more than a dozen Telegram accounts and identified the phone numbers of 15 million Iranian users, as well as the associated user IDs. Researchers said the hackers belonged to a group known as Rocket Kitten. Rocket Kitten's attacks were similar to ones attributed to Iran's Islamic Revolutionary Guards Corps. The attackers took advantage of a programming interface built into Telegram. According to Telegram, these mass checks are no longer possible because of limitations introduced into its API earlier in 2016.
Login SMS messages are known to have been intercepted in Iran, Russia and Germany, possibly in coordination with phone or telecom companies. Pavel Durov has said that Telegram users in "troubled countries" should enable two-factor authentication by creating passwords in order to prevent this.
In June 2017, Pavel Durov in an interview said that U.S. intelligence agencies tried to bribe the company's developers to weaken Telegram's encryption or install a backdoor during their visit to the U.S. in 2016.
In 2018, Telegram sent a message to all Iranian users stating that the Telegram Talai and Hotgram unofficial clients are not secure.
In March 2014, Telegram promised that "all code will be released eventually", including all the various client applications (Android, iOS, desktop, etc.) and the server-side code. As of May 2021, Telegram had not published their server-side source code. In January 2021, Durov explained his rationale for not releasing server-side code, citing reasons such as inability for end-users to verify that the released code is the same code run on servers, and a government that wanted to acquire the server code and make an instant messaging network that would end competitors.
On 9 June 2019, The Intercept released leaked Telegram messages exchanged between current Brazilian Minister of Justice and former judge Sérgio Moro and federal prosecutors. The hypothesis is that either mobile devices were hacked by SIM swap or the targets' computers were compromised. The Telegram team tweeted that it was either because the user had malware or they were not using two-step verification.
On 12 June 2019, Telegram confirmed that it suffered a denial-of-service attack which disrupted normal app functionality for approximately one hour. Pavel Durov tweeted that the IP addresses used in the attack mostly came from China.
In December 2019, multiple Russian businessmen suffered account takeovers that involved bypassing SMS single-factor authentication. Security company Group-IB suggested SS7 mobile signalling protocol weaknesses, illegal usage of surveillance equipment, or telecom insider attacks.
On 30 March 2020, an Elasticsearch database holding 42 million records containing user IDs and phone numbers of Iranian users was exposed online without a password. The accounts were extracted from not Telegram but an unofficial version of Telegram, in what appears to be a possibly government-sanctioned fork. It took 11 days for the database to be taken down, but the researchers say the data was accessed by other parties, including a hacker who reported the information to a specialized forum.
In September 2020, it was reported that Iran's RampantKitten espionage group ran a phishing and surveillance campaign against dissidents on Telegram. The attack relied on people downloading a malware-infected file from any source, at which point it would replace Telegram files on the device and 'clone' session data. David Wolpoff, a former Department of Defense contractor, has stated that the weak link in the attack was the device itself and not any of the affected apps: "There's no way for a secure communication app to keep a user safe when the end devices are compromised."
In July 2021, researchers from Royal Holloway, University of London and ETH Zurich published an analysis of the MTProto protocol, concluding that the protocol could provide a "confidential and integrity-protected channel" for communication. They also found that attackers had the theoretical ability to reorder messages coming from the client to the server though the attacker would not be able to see the content of the messages. Several other theoretical vulnerabilities were reported as well, in response to which Telegram released a document stating that the MITM attack on the key exchange was impossible as well as detailing the changes made to the protocol to protect from it in the future. All issues were patched before the paper's publication with a security bounty paid out to the researchers.
In September 2021, a Russian researcher published details about a bug with the self-destruct feature that allowed the user to recover deleted photos from their own device. The bug was patched prior to publication and Telegram representatives offered a €1,000 bug bounty. The researcher did not sign the NDA that came with the offer and did not receive the award, opting to disclose the bug.
In March 2023, the Norwegian National Security Authority (NSM) advised against the use of Telegram and TikTok on business devices (especially the ones used for government related activities), the assessment has been commissioned and supported by the Ministry of Justice and Public Security, Emilie Enger Mehl. Regarding Telegram, the report cites its lack of end-to-end encryption by default, its Russian origins and third-party open source intelligence as major critical points.
In July 2024, ESET reported a vulnerability allowed malicious files being sent to users masked in multimedia.
Cryptography contests
Telegram has organized two cryptography contests to challenge its own security. Third parties were asked to break the service's cryptography and disclose the information contained within a secret chat between two computer-controlled users. A reward of respectively and was offered. Both of these contests expired with no winners.
Security researcher Moxie Marlinspike, founder of the competing Signal messenger, and commenters on Hacker News criticized the first contest for being rigged or framed in Telegram's favor and said that Telegram's statements on the value of these contests as proof of the cryptography's quality are misleading. This was because the cryptography contest could not be won even with completely broken algorithms such as MD2 (hash function) used as key stream extractor, and primitives such as the Dual EC DRBG that is known to be backdoored.
Censorship
Telegram has been blocked temporarily or permanently by some governments including Iran, China, Brazil, and Pakistan. The Russian government blocked Telegram for several years before lifting the ban in 2020. The company's founder has said he wants the app to have an anti-censorship tool for Iran and China similar to the app's role in fighting censorship in Russia. On 19 April 2024, Apple removed Telegram from the App Store in China.
In September 2024, Ukraine banned the usage of Telegram by government officials, military personnel, and key workers on official devices, citing fears of Russian espionage. The National Security and Defense Council of Ukraine enforced the restrictions following evidence reported to them by Kyrylo Budanov, showing Russia's ability to access messages and user data on the platform. Andriy Kovalenko, head of the security council's centre on countering disinformation, clarified that the ban was limited to official devices and did not extend to personal phones.
2019 Puerto Rico "Telegramgate"
Telegram was the main subject surrounding the 2019 Puerto Rico riots that ended up in the resignation of then-Governor Ricardo Rosselló. Hundreds of pages of a group chat between Rosselló and members of his staff were leaked. The messages were considered vulgar, racist, and homophobic, with members of the chat discussing how they would use the media to target potential political opponents.
2021 shutdown of Russian political bots
In September 2021, prior to the regional elections in Russia, Telegram suspended several bots spreading information about the election, including a bot run by the opposition party and critics of incumbent president Vladimir Putin's government, citing election silence as the reason, though a blog post by the company's CEO implied the company was following Apple and Google, which "dictate the rules of the game to developers". The blocking of the main Smart Voting bot was criticized by allies of Alexei Navalny, a Kremlin critic and former opposition leader. Navalny's spokeswoman Kira Yarmysh called the block and the deletion of the tactical voting app from app stores "censorship [...] imposed by private companies".
In a later blog post, Durov directly stated that the block was a result of pressure from Google and Apple as refusal to comply with their policies would result "in an immediate shutdown of Telegram for millions of users". The post included a screenshot showing an internal email sent by the App Store to developers, demanding the takedown of content related to Navalny.
2022 Delhi High Court ruling
On 24 November 2022, Telegram disclosed the admin names, phone numbers and IP addresses of channels accused of unauthorised sharing of national exam study materials following an order by the Delhi High Court which rejected Telegram's argument that its regional servers were located in Singapore and thus no data could be disclosed as the local laws prohibit it.
2024 arrest of Pavel Durov
On 24 August 2024, Pavel Durov, who is both a French and UAE citizen, was arrested in France by French authorities and four days later charged with a wide array of crimes, including complicity in managing an online platform to enable illegal transactions; complicity in crimes such as enabling the distribution of child sexual abuse material; drug trafficking and fraud; and a refusal to cooperate with law enforcement authorities. Durov posted bail of five million Euros, was barred from leaving France, and was released on condition he report to a French police station twice weekly. The case would be handled by a special magistrate with investigative and prosecutorial powers.
| Technology | Social network and blogging | null |
29188721 | https://en.wikipedia.org/wiki/RNA%20interference | RNA interference | RNA interference (RNAi) is a biological process in which RNA molecules are involved in sequence-specific suppression of gene expression by double-stranded RNA, through translational or transcriptional repression. Historically, RNAi was known by other names, including co-suppression, post-transcriptional gene silencing (PTGS), and quelling. The detailed study of each of these seemingly different processes elucidated that the identity of these phenomena were all actually RNAi. Andrew Fire and Craig Mello shared the 2006 Nobel Prize in Physiology or Medicine for their work on RNAi in the nematode worm Caenorhabditis elegans, which they published in 1998. Since the discovery of RNAi and its regulatory potentials, it has become evident that RNAi has immense potential in suppression of desired genes. RNAi is now known as precise, efficient, stable and better than antisense therapy for gene suppression. Antisense RNA produced intracellularly by an expression vector may be developed and find utility as novel therapeutic agents.
Two types of small ribonucleic acid (RNA) molecules, microRNA (miRNA) and small interfering RNA (siRNA), are central to components to the RNAi pathway. Once mRNA is degraded, post-transcriptional silencing occurs as protein translation is prevented. Transcription can be inhibited via the pre-transcriptional silencing mechanism of RNAi, through which an enzyme complex catalyzes DNA methylation at genomic positions complementary to complexed siRNA or miRNA. RNAi has an important role in defending cells against parasitic nucleotide sequences (e.g., viruses or transposons) and also influences development of organisms.
The RNAi pathway is a naturally occurring process found in many eukaryotes and animal cells. It is initiated by the enzyme Dicer, which cleaves long double-stranded RNA (dsRNA) molecules into short double-stranded fragments of approximately 21 to 23 nucleotide siRNAs. Each siRNA is unwound into two single-stranded RNAs (ssRNAs), the passenger (sense) strand and the guide (antisense) strand. The passenger strand is then cleaved by the protein Argonaute 2 (Ago2). The passenger strand is degraded and the guide strand is incorporated into the RNA-induced silencing complex (RISC). The RISC assembly then binds and degrades the target mRNA. Specifically, this is accomplished when the guide strand pairs with a complementary sequence in a mRNA molecule and induces cleavage by Ago2, a catalytic component of the RISC. In some organisms, this process spreads systemically, despite the initially limited molar concentrations of siRNA.
RNAi is a valuable research tool, both in cell culture and in living organisms, because synthetic dsRNA introduced into cells can selectively and robustly induce suppression of specific genes of interest. RNAi may be used for large-scale screens that systematically shut down each gene (and the subsequent proteins it codes for) in the cell, which can help to identify the components necessary for a particular cellular process or an event such as cell division. The pathway is also used as a practical tool for food, medicine and insecticides.
Cellular mechanism
RNAi is an RNA-dependent gene silencing process that is controlled by RISC and is initiated by short double-stranded RNA molecules in a cell's cytoplasm, where they interact with the catalytic RISC component Argonaute. When the dsRNA is exogenous (coming from infection by a virus with an RNA genome or laboratory manipulations), the RNA is imported directly into the cytoplasm and cleaved to short fragments by Dicer. The initiating dsRNA can also be endogenous (originating in the cell), as in pre-microRNAs expressed from RNA-coding genes in the genome. The primary transcripts from such genes are first processed to form the characteristic stem-loop structure of pre-miRNA in the nucleus, then exported to the cytoplasm. Thus, the two dsRNA pathways, exogenous and endogenous, converge at the RISC.
Exogenous dsRNA initiates RNAi by activating the ribonuclease protein Dicer, which binds and cleaves dsRNAs in plants, or short hairpin RNAs (shRNAs) in humans, to produce double-stranded fragments of 20–25 base pairs with a 2-nucleotide overhang at the 3′ end. Bioinformatics studies on the genomes of multiple organisms suggest this length maximizes target-gene specificity and minimizes non-specific effects. These short double-stranded fragments are called siRNAs. These siRNAs are then separated into single strands and integrated into an active RISC, by RISC-Loading Complex (RLC). RLC includes Dicer-2 and R2D2, and is crucial to unite Ago2 and RISC. TATA-binding protein-associated factor 11 (TAF11) assembles the RLC by facilitating Dcr-2-R2D2 tetramerization, which increases the binding affinity to siRNA by 10-fold. Association with TAF11 would convert the R2-D2-Initiator (RDI) complex into the RLC. R2D2 carries tandem double-stranded RNA-binding domains to recognize the thermodynamically stable terminus of siRNA duplexes, whereas Dicer-2 the other less stable extremity. Loading is asymmetric: the MID domain of Ago2 recognizes the thermodynamically stable end of the siRNA. Therefore, the "passenger" (sense) strand whose 5′ end is discarded by MID is ejected, while the saved "guide" (antisense) strand cooperates with AGO to form the RISC.
After integration into the RISC, siRNAs base-pair to their target mRNA and cleave it, thereby preventing it from being used as a translation template. Differently from siRNA, a miRNA-loaded RISC complex scans cytoplasmic mRNAs for potential complementarity. Instead of destructive cleavage (by Ago2), miRNAs rather target the 3′ untranslated region (UTR) regions of mRNAs where they typically bind with imperfect complementarity, thus blocking the access of ribosomes for translation.
Exogenous dsRNA is detected and bound by an effector protein, known as RDE-4 in C. elegans and R2D2 in Drosophila, that stimulates Dicer activity. The mechanism producing this length specificity is unknown and this protein only binds long dsRNAs.
In C. elegans this initiation response is amplified through the synthesis of a population of 'secondary' siRNAs during which the Dicer-produced initiating or 'primary' siRNAs are used as templates. These 'secondary' siRNAs are structurally distinct from Dicer-produced siRNAs and appear to be produced by an RNA-dependent RNA polymerase (RdRP).
MicroRNA
MicroRNAs (miRNAs) are genomically encoded non-coding RNAs that help regulate gene expression, particularly during development. The phenomenon of RNAi, broadly defined, includes the endogenously induced gene silencing effects of miRNAs as well as silencing triggered by foreign dsRNA. Mature miRNAs are structurally similar to siRNAs produced from exogenous dsRNA, but before reaching maturity, miRNAs must first undergo extensive post-transcriptional modification. A miRNA is expressed from a much longer RNA-coding gene as a primary transcript known as a pri-miRNA which is processed, in the cell nucleus, to a 70-nucleotide stem-loop structure called a pre-miRNA by the microprocessor complex. This complex consists of an RNase III enzyme called Drosha and a dsRNA-binding protein DGCR8. The dsRNA portion of this pre-miRNA is bound and cleaved by Dicer to produce the mature miRNA molecule that can be integrated into the RISC complex; thus, miRNA and siRNA share the same downstream cellular machinery. First, viral encoded miRNA was described in Epstein–Barr virus (EBV). Thereafter, an increasing number of microRNAs have been described in viruses. VIRmiRNA is a comprehensive catalogue covering viral microRNA, their targets and anti-viral miRNAs (see also VIRmiRNA resource: http://crdd.osdd.net/servers/virmirna/).
siRNAs derived from long dsRNA precursors differ from miRNAs in that miRNAs, especially those in animals, typically have incomplete base pairing to a target and inhibit the translation of many different mRNAs with similar sequences. In contrast, siRNAs typically base-pair perfectly and induce mRNA cleavage only in a single, specific target. In Drosophila and C. elegans, miRNA and siRNA are processed by distinct Argonaute proteins and Dicer enzymes.
Three prime untranslated regions and microRNAs
Three prime untranslated regions (3′UTRs) of mRNAs often contain regulatory sequences that post-transcriptionally cause RNAi. Such 3′-UTRs often contain both binding sites for miRNAs as well as for regulatory proteins. By binding to specific sites within the 3′-UTR, miRNAs can decrease gene expression of various mRNAs by either inhibiting translation or directly causing degradation of the transcript. The 3′-UTR also may have silencer regions that bind repressor proteins that inhibit the expression of a mRNA.
The 3′-UTR often contains microRNA response elements (MREs). MREs are sequences to which miRNAs bind. These are prevalent motifs within 3′-UTRs. Among all regulatory motifs within the 3′-UTRs (e.g. including silencer regions), MREs make up about half of the motifs.
As of 2023, the miRBase web site, an archive of miRNA sequences and annotations, listed 28,645 entries in 271 biologic species. Of these, 1,917 miRNAs were in annotated human miRNA loci. miRNAs were predicted to have an average of about four hundred target mRNAs (affecting expression of several hundred genes). Friedman et al. estimate that >45,000 miRNA target sites within human mRNA 3′UTRs are conserved above background levels, and >60% of human protein-coding genes have been under selective pressure to maintain pairing to miRNAs.
Direct experiments show that a single miRNA can reduce the stability of hundreds of unique mRNAs. Other experiments show that a single miRNA may repress the production of hundreds of proteins, but that this repression often is relatively mild (less than 2-fold).
The effects of miRNA dysregulation of gene expression seem to be important in cancer. For instance, in gastrointestinal cancers, nine miRNAs have been identified as epigenetically altered and effective in down regulating DNA repair enzymes.
The effects of miRNA dysregulation of gene expression also seem to be important in neuropsychiatric disorders, such as schizophrenia, bipolar disorder, major depression, Parkinson's disease, Alzheimer's disease and autism spectrum disorders.
RISC activation and catalysis
Exogenous dsRNA is detected and bound by an effector protein, known as RDE-4 in C. elegans and R2D2 in Drosophila, that stimulates Dicer activity. This protein only binds long dsRNAs, but the mechanism producing this length specificity is unknown. This RNA-binding protein then facilitates the transfer of cleaved siRNAs to the RISC complex.
In C. elegans this initiation response is amplified through the synthesis of a population of 'secondary' siRNAs during which the Dicer-produced initiating or 'primary' siRNAs are used as templates. These 'secondary' siRNAs are structurally distinct from Dicer-produced siRNAs and appear to be produced by an RNA-dependent RNA polymerase (RdRP).
The active components of an RNA-induced silencing complex (RISC) are endonucleases called Argonaute proteins, which cleave the target mRNA strand complementary to their bound siRNA. As the fragments produced by Dicer are double-stranded, they could each in theory produce a functional siRNA. However, only one of the two strands, which is known as the guide strand, binds Argonaute and directs gene silencing. The other anti-guide strand or passenger strand is degraded during RISC activation. Although it was first believed that an ATP-dependent helicase separated these two strands, the process proved to be ATP-independent and performed directly by the protein components of RISC. However, an in vitro kinetic analysis of RNAi in the presence and absence of ATP showed that ATP may be required to unwind and remove the cleaved mRNA strand from the RISC complex after catalysis. The guide strand tends to be the one whose 5′ end is less stably paired to its complement, but strand selection is unaffected by the direction in which Dicer cleaves the dsRNA before RISC incorporation. Instead, the R2D2 protein may serve as the differentiating factor by binding the more-stable 5′ end of the passenger strand.
The structural basis for binding of RNA to the Argonaute protein was examined by X-ray crystallography of the binding domain of an RNA-bound Argonaute. Here, the phosphorylated 5′ end of the RNA strand enters a conserved basic surface pocket and makes contacts through a divalent cation (an atom with two positive charges) such as magnesium and by aromatic stacking (a process that allows more than one atom to share an electron by passing it back and forth) between the 5′ nucleotide in the siRNA and a conserved tyrosine residue. This site is thought to form a nucleation site for the binding of the siRNA to its mRNA target. Analysis of the inhibitory effect of mismatches in either the 5’ or 3’ end of the guide strand has demonstrated that the 5’ end of the guide strand is likely responsible for matching and binding the target mRNA, while the 3’ end is responsible for physically arranging target mRNA into a cleavage-favorable RISC region.
It is not understood how the activated RISC complex locates complementary mRNAs within the cell. Although the cleavage process has been proposed to be linked to translation, translation of the mRNA target is not essential for RNAi-mediated degradation. Indeed, RNAi may be more effective against mRNA targets that are not translated. Argonaute proteins are localized to specific regions in the cytoplasm called P-bodies (also cytoplasmic bodies or GW bodies), which are regions with high rates of mRNA decay; miRNA activity is also clustered in P-bodies. Disruption of P-bodies decreases the efficiency of RNAi, suggesting that they are a critical site in the RNAi process.
Transcriptional silencing
Components of the RNAi pathway are used in many eukaryotes in the maintenance of the organization and structure of their genomes. Modification of histones and associated induction of heterochromatin formation serves to downregulate genes pre-transcriptionally; this process is referred to as RNA-induced transcriptional silencing (RITS), and is carried out by a complex of proteins called the RITS complex. In fission yeast this complex contains Argonaute, a chromodomain protein Chp1, and a protein called Tas3 of unknown function. As a consequence, the induction and spread of heterochromatic regions requires the Argonaute and RdRP proteins. Indeed, deletion of these genes in the fission yeast S. pombe disrupts histone methylation and centromere formation, causing slow or stalled anaphase during cell division. In some cases, similar processes associated with histone modification have been observed to transcriptionally upregulate genes.
The mechanism by which the RITS complex induces heterochromatin formation and organization is not well understood. Most studies have focused on the mating-type region in fission yeast, which may not be representative of activities in other genomic regions/organisms. In maintenance of existing heterochromatin regions, RITS forms a complex with siRNAs complementary to the local genes and stably binds local methylated histones, acting co-transcriptionally to degrade any nascent pre-mRNA transcripts that are initiated by RNA polymerase. The formation of such a heterochromatin region, though not its maintenance, is Dicer-dependent, presumably because Dicer is required to generate the initial complement of siRNAs that target subsequent transcripts. Heterochromatin maintenance has been suggested to function as a self-reinforcing feedback loop, as new siRNAs are formed from the occasional nascent transcripts by RdRP for incorporation into local RITS complexes. The relevance of observations from fission yeast mating-type regions and centromeres to mammals is not clear, as heterochromatin maintenance in mammalian cells may be independent of the components of the RNAi pathway.
Crosstalk with RNA editing
The type of RNA editing that is most prevalent in higher eukaryotes converts adenosine nucleotides into inosine in dsRNAs via the enzyme adenosine deaminase (ADAR). It was originally proposed in 2000 that the RNAi and A→I RNA editing pathways might compete for a common dsRNA substrate. Some pre-miRNAs do undergo A→I RNA editing and this mechanism may regulate the processing and expression of mature miRNAs. Furthermore, at least one mammalian ADAR can sequester siRNAs from RNAi pathway components. Further support for this model comes from studies on ADAR-null C. elegans strains indicating that A→I RNA editing may counteract RNAi silencing of endogenous genes and transgenes.
Variation among organisms
Organisms vary in their ability to take up foreign dsRNA and use it in the RNAi pathway. The effects of RNAi can be both systemic and heritable in plants and C. elegans, although not in Drosophila or mammals. In plants, RNAi is thought to propagate by the transfer of siRNAs between cells through plasmodesmata (channels in the cell walls that enable communication and transport). Heritability comes from methylation of promoters targeted by RNAi; the new methylation pattern is copied in each new generation of the cell. A broad general distinction between plants and animals lies in the targeting of endogenously produced miRNAs; in plants, miRNAs are usually perfectly or nearly perfectly complementary to their target genes and induce direct mRNA cleavage by RISC, while animals' miRNAs tend to be more divergent in sequence and induce translational repression. This translational effect may be produced by inhibiting the interactions of translation initiation factors with the mRNA's polyadenine tail.
Some eukaryotic protozoa such as Leishmania major and Trypanosoma cruzi lack the RNAi pathway entirely. Most or all of the components are also missing in some fungi, most notably the model organism Saccharomyces cerevisiae. The presence of RNAi in other budding yeast species such as Saccharomyces castellii and Candida albicans, further demonstrates that inducing two RNAi-related proteins from S. castellii facilitates RNAi in S. cerevisiae. That certain ascomycetes and basidiomycetes are missing RNAi pathways indicates that proteins required for RNA silencing have been lost independently from many fungal lineages, possibly due to the evolution of a novel pathway with similar function, or to the lack of selective advantage in certain niches.
Related prokaryotic systems
Gene expression in prokaryotes is influenced by an RNA-based system similar in some respects to RNAi. Here, RNA-encoding genes control mRNA abundance or translation by producing a complementary RNA that anneals to an mRNA. However these regulatory RNAs are not generally considered to be analogous to miRNAs because the Dicer enzyme is not involved. It has been suggested that CRISPR interference systems in prokaryotes are analogous to eukaryotic RNAi systems, although none of the protein components are orthologous.
Biological functions
Immunity
RNAi is a vital part of the immune response to viruses and other foreign genetic material, especially in plants where it may also prevent the self-propagation of transposons. Plants such as Arabidopsis thaliana express multiple Dicer homologs that are specialized to react differently when the plant is exposed to different viruses. Even before the RNAi pathway was fully understood, it was known that induced gene silencing in plants could spread throughout the plant in a systemic effect and could be transferred from stock to scion plants via grafting. This phenomenon has since been recognized as a feature of the plant immune system which allows the entire plant to respond to a virus after an initial localized encounter. In response, many plant viruses have evolved elaborate mechanisms to suppress the RNAi response. These include viral proteins that bind short double-stranded RNA fragments with single-stranded overhang ends, such as those produced by Dicer. Some plant genomes also express endogenous siRNAs in response to infection by specific types of bacteria. These effects may be part of a generalized response to pathogens that downregulates any metabolic process in the host that aids the infection process.
Although animals generally express fewer variants of the Dicer enzyme than plants, RNAi in some animals produces an antiviral response. In both juvenile and adult Drosophila, RNAi is important in antiviral innate immunity and is active against pathogens such as Drosophila X virus. A similar role in immunity may operate in C. elegans, as Argonaute proteins are upregulated in response to viruses and worms that overexpress components of the RNAi pathway are resistant to viral infection.
The role of RNAi in mammalian innate immunity is poorly understood, and relatively little data is available. However, the existence of viruses that encode genes able to suppress the RNAi response in mammalian cells may be evidence in favour of an RNAi-dependent mammalian immune response, although this hypothesis has been challenged as poorly substantiated.
Evidence for the existence of a functional antiviral RNAi pathway in mammalian cells has been presented.
Other functions for RNAi in mammalian viruses also exist, such as miRNAs expressed by the herpes virus that may act as heterochromatin organization triggers to mediate viral latency.
Downregulation of genes
Endogenously expressed miRNAs, including both intronic and intergenic miRNAs, are most important in translational repression and in the regulation of development, especially on the timing of morphogenesis and the maintenance of undifferentiated or incompletely differentiated cell types such as stem cells. The role of endogenously expressed miRNA in downregulating gene expression was first described in C. elegans in 1993. In plants this function was discovered when the "JAW microRNA" of Arabidopsis was shown to be involved in the regulation of several genes that control plant shape. In plants, the majority of genes regulated by miRNAs are transcription factors; thus miRNA activity is particularly wide-ranging and regulates entire gene networks during development by modulating the expression of key regulatory genes, including transcription factors as well as F-box proteins. In many organisms, including humans, miRNAs are linked to the formation of tumors and dysregulation of the cell cycle. Here, miRNAs can function as both oncogenes and tumor suppressors.
Evolution
Based on parsimony-based phylogenetic analysis, the most recent common ancestor of all eukaryotes most likely already possessed an early RNAi pathway; the absence of the pathway in certain eukaryotes is thought to be a derived characteristic. This ancestral RNAi system probably contained at least one Dicer-like protein, one Argonaute, one PIWI protein, and an RNA-dependent RNA polymerase that may also have played other cellular roles. A large-scale comparative genomics study likewise indicates that the eukaryotic crown group already possessed these components, which may then have had closer functional associations with generalized RNA degradation systems such as the exosome. This study also suggests that the RNA-binding Argonaute protein family, which is shared among eukaryotes, most archaea, and at least some bacteria (such as Aquifex aeolicus), is homologous to and originally evolved from components of the translation initiation system.
Applications
RNAi pathway for gene knockdown
Gene knockdown is a method used to reduce the expression of an organism’s specific genes. This is accomplished by using the naturally occurring process of RNAi. This gene knockdown technique uses a double-stranded siRNA molecule that is synthesized with a sequence complementary to the gene of interest. The RNAi cascade begins once the Dicer enzyme starts to process siRNA. The end result of the process leads to degradation of mRNA and destroys any instructions needed to build certain proteins. Using this method, researchers are able to decrease (but not completely eliminate) the expression of a targeted gene. Studying the effects of this decrease in expression may show the physiological role or impact of the targeted gene products.
Off-Target Effects of Gene Knockdown
Extensive efforts in computational biology have been directed toward the design of successful dsRNA reagents that maximize gene knockdown but minimize "off-target" effects. Off-target effects arise when an introduced RNA has a base sequence that can pair with and thus reduce the expression of multiple genes. Such problems occur more frequently when the dsRNA contains repetitive sequences. It has been estimated from studying the genomes of humans, C. elegans and S. pombe that about 10% of possible siRNAs have substantial off-target effects. A multitude of software tools have been developed implementing algorithms for the design of general mammal-specific, and virus-specific siRNAs that are automatically checked for possible cross-reactivity.
Depending on the organism and experimental system, the exogenous RNA may be a long strand designed to be cleaved by Dicer, or short RNAs designed to serve as siRNA substrates. In most mammalian cells, shorter RNAs are used because long double-stranded RNA molecules induce the mammalian interferon response, a form of innate immunity that reacts nonspecifically to foreign genetic material. Mouse oocytes and cells from early mouse embryos lack this reaction to exogenous dsRNA and are therefore a common model system for studying mammalian gene-knockdown effects. Specialized laboratory techniques have also been developed to improve the utility of RNAi in mammalian systems by avoiding the direct introduction of siRNA, for example, by stable transfection with a plasmid encoding the appropriate sequence from which siRNAs can be transcribed, or by more elaborate lentiviral vector systems allowing the inducible activation or deactivation of transcription, known as conditional RNAi.
Medications
The technique of knocking down genes using RNAi therapeutics has demonstrated success in randomized controlled clinical studies. These medications are a growing class of siRNA-based drugs that decrease the expression of proteins encoded by certain genes. To date, five RNAi medications have been approved by regulatory authorities in the US and Europe: patisiran (2018), givosiran (2019), lumasiran (2020), inclisiran (2020 in Europe with anticipated US approval in 2021), and vutrisiran (2022).
While all of the current regulatory body approved RNAi therapeutics focus on diseases that originate in the liver, additional medications under investigation target a host of disease areas including cardiovascular diseases, bleeding disorders, alcohol use disorders, cystic fibrosis, gout, carcinoma, and eye disorders.
Patisiran is the first double stranded siRNA-based medication approved in 2018 and developed by Alnylam Pharmaceuticals. Patisiran uses the RNAi cascade to suppress the gene that codes for TTR (transthryetin). Mutations in this gene may cause the misfolding of a protein responsible for hereditary ATTR amyloidosis. To achieve therapeutic response, patisiran is encased by a lipid nanoparticle membrane that facilitates crossover into the cytoplasm. Once inside the cell, the siRNA begins processing by the enzyme Dicer. Patisiran is administered by a healthcare professional through an intravenous infusion with dosing based on body weight. Warnings and precautions include risk of infusion-related reactions and reduced vitamin A levels (serum).
In 2019, the FDA and EMA approved givosiran for the treatment of adults with acute hepatic porphyria (AHP). The FDA also granted givosiran a breakthrough therapy designation, priority review designation, and orphan drug designation for the treatment of acute hepatic porphyria (AHP) in November 2019. By 2020, givosiran received EMA approval. Givosiran is an siRNA that breaks down aminolevulinic acid synthase 1 (ALAS1) mRNA in the liver. Breaking down ALAS1 mRNA prevents toxins (responsible for neurovisceral attacks and AHP disease) such as aminolevulinic acid (ALA) and porphobilinogen (PBG) from accumulating. To facilitate entry into the cytoplasm, givosiran uses GalNAc ligands and enters into liver cells. The medication is administered subcutaneously by a healthcare professional with dosing based on body weight. Warnings and precautions include risk of anaphylactic reactions, hepatic toxicity, renal toxicity and injection site reactions.
Lumasiran was approved as a siRNA-based medication in 2020 for use in both the European Union and the United States. This medication is used for the treatment of primary hyperoxaluria type 1 (PH1) in pediatric and adult populations. The drug is designed to reduce hepatic oxalate production and urinary oxalate levels through RNAi by targeting hydroxyacid oxidase 1 (HAO1) mRNA for breakdown. Lowering HAO1 enzyme levels reduces the oxidation of glycolate to glyoxylate (which is a substrate for oxalate). Lumasiran is administered subcutaneously by a healthcare professional with dosing based on body weight. Data from randomized controlled clinical trials indicate that the most common adverse reaction that was reported was injection site reactions. These reactions were mild and were present in 38 percent of patients treated with lumasiran.
In 2022, the FDA and EMA approved vutrisiran for the treatment of adults with hereditary transthyretin mediated amyloidosis with polyneuropathy stage 1 or 2. Vutrisiran is designed to break down the mRNA that codes for transthyretin.
Other investigational drugs using RNAi that are being developed by pharmaceutical companies such as Arrowhead Pharmaceuticals, Dicerna, Alnylam Pharmaceuticals, Amgen, and Sylentis. These medications cover a variety of targets via RNAi and diseases.
Investigational RNAi therapeutics in development:
Legal categorization and legal issues in a near future
Currently, both miRNA and SiRNA are currently chemically synthesized and so, are legally categorized inside EU and in USA as "simple" medicinal products. But as bioengineered siRNA (BERAs) are in development, these would be classified as biological medicinal products, at least in EU. The development of the BERAs technology raises the question of the categorization of drugs having the same mechanism of action but being produced chemically or biologically. This lack of consistency should be addressed.
Delivery mechanisms
To achieve the clinical potential of RNAi, siRNA must be efficiently transported to the cells of target tissues. However, there are various barriers that must be fixed before it can be used clinically. For example, "naked" siRNA is susceptible to several obstacles that reduce its therapeutic efficacy. Additionally, once siRNA has entered the bloodstream, naked RNA can be degraded by serum nucleases and can stimulate the innate immune system. Due to its size and highly polyanionic (containing negative charges at several sites) nature, unmodified siRNA molecules cannot readily enter the cells through the cell membrane. Therefore, artificial or nanoparticle encapsulated siRNA must be used. If siRNA is transferred across the cell membrane, unintended toxicities can occur if therapeutic doses are not optimized, and siRNAs can exhibit off-target effects (e.g. unintended downregulation of genes with partial sequence complementarity). Even after entering the cells, repeated dosing is required since their effects are diluted at each cell division. In response to these potential issues and barriers, two approaches help facilitate siRNA delivery to target cells: lipid nanoparticles and conjugates.
Lipid nanoparticles
Lipid nanoparticles (LNPs) are based on liposome-like structures that are typically made of an aqueous center surrounded by a lipid shell. A subset of liposomal structures used for delivery drugs to tissues rest in large unilamellar vesicles (LUVs) which may be 100 nm in size. LNP delivery mechanisms have become an increasing source of encasing nucleic acids and may include plasmids, CRISPR and mRNA.
The first approved use of lipid nanoparticles as a drug delivery mechanism began in 2018 with the siRNA drug patisiran, developed by Alnylam Pharmaceuticals. Dicerna Pharmaceuticals, Persomics, Sanofi and Sirna Therapeutics also worked to bring RNAi therapies to market.
Other recent applications include two FDA approved COVID-19 vaccines: mRNA-1273, developed by Moderna. and BNT162b, developed by a collaboration between Pfizer and BioNtech. These two vaccines use lipid nanoparticles to deliver antigen mRNA. Encapsulating the mRNA molecule in lipid nanoparticles was a critical breakthrough for producing viable mRNA vaccines, solving a number of key technical barriers in delivering the mRNA molecule into the host cell as distributed through apolipoprotein E (apoE) in the low-density lipoprotein receptor (LDLR). In December 2020, Novartis announced that positive results from phase III efficacy studies deemed inclisiran was a treatment for heterozygous familial hypercholesterolemia (HeFH) and atherosclerotic cardiovascular disease (ASCVD).
Conjugates
In addition to LNPs, RNAi therapeutics have targeted delivery through siRNA conjugates (e.g., GalNAc, carbohydrates, peptides, aptamers, antibodies). Therapeutics using siRNA conjugates have been developed for rare or genetic diseases such as acute hepatic porphyria (AHP), hemophilia, primary hyperoxaluria (PH) and hereditary ATTR amyloidosis as well as other cardiometabolic diseases such as hypertension and non-alcoholic steatohepatitis (NASH).
Biotechnology
RNAi has been used for a variety of other applications including food, crops and insecticides. The use of the RNAi pathway has developed numerous products such as foods like Arctic apples, nicotine-free tobacco, decaffeinated coffee, nutrient fortified vegetation and hypoallergenic crops. The emerging use of RNAi has the potential to develop many other products for future use.
Viral infection
Antiviral treatment is one of the earliest proposed RNAi-based medical applications, and two different types have been developed. The first type is to target viral RNAs. Many studies have shown that targeting viral RNAs can suppress the replication of numerous viruses, including HIV, HPV, hepatitis A, hepatitis B, influenza virus, respiratory syncytial virus (RSV), SARS coronavirus (SARS-CoV), adenovirus and measles virus. The other strategy is to block the initial viral entries by targeting the host cell genes. For example, suppression of chemokine receptors (CXCR4 and CCR5) on host cells can prevent HIV viral entry.
Cancer
While traditional chemotherapy can effectively kill cancer cells, lack of specificity for discriminating normal cells and cancer cells in these treatments usually cause severe side effects. Numerous studies have demonstrated that RNAi can provide a more specific approach to inhibit tumor growth by targeting cancer-related genes (i.e., oncogene). It has also been proposed that RNAi can enhance the sensitivity of cancer cells to chemotherapeutic agents, providing a combinatorial therapeutic approach with chemotherapy. Another potential RNAi-based treatment is to inhibit cell invasion and migration.
Compared with chemotherapy or other anti-cancer drugs, there are a lot of advantages of siRNA drug. SiRNA acts on the post-transcriptional stage of gene expression, so it does not modify or change DNA in a deleterious effect. SiRNA can also be used to produce a specific response in a certain type of way, such as by downgrading suppression of gene expression. In a single cancer cell, siRNA can cause dramatic suppression of gene expression with just several copies. This happens by silencing cancer-promoting genes with RNAi, as well as targeting an mRNA sequence.
RNAi drugs treat cancer by silencing certain cancer promoting genes. This is done by complementing the cancer genes with the RNAi, such as keeping the mRNA sequences in accordance with the RNAi drug. Ideally, RNAi is should be injected and/or chemically modified so the RNAi can reach cancer cells more efficiently. RNAi uptake and regulation is controlled by the kidneys.
Neurological diseases
RNAi strategies also show potential for treating neurodegenerative diseases. Studies in cells and in mouse have shown that specifically targeting Amyloid beta-producing genes (e.g. BACE1 and APP) by RNAi can significantly reduced the amount of Aβ peptide which is correlated with the cause of Alzheimer's disease. In addition, this silencing-based approaches also provide promising results in treatment of Parkinson's disease and Polyglutamine disease.
Stimulation of immune response
The human immune system is divided into two separate branches: the innate immune system and the adaptive immune system. The innate immune system is the first defense against infection and responds to pathogens in a generic fashion. On the other hand, the adaptive immune system, a system that was evolved later than the innate, is composed mainly of highly specialized B and T cells that are trained to react to specific portions of pathogenic molecules.
The challenge between old pathogens and new has helped create a system of guarded cells and particles that are called safe framework. This framework has given humans an army of systems that search out and destroy invader particles, such as pathogens, microscopic organisms, parasites, and infections. The mammalian safe framework has developed to incorporate siRNA as a tool to indicate viral contamination, which has allowed siRNA is create an intense innate immune response.
siRNA is controlled by the innate immune system, which can be divided into the acute inflammatory responses and antiviral responses. The inflammatory response is created with signals from small signaling molecules, or cytokines. These include interleukin-1 (IL-1), interleukin-6 (IL-6), interleukin-12 (IL-12) and tumor necrosis factor α (TNF-α). The innate immune system generates inflammation and antiviral responses, which cause the release pattern recognition receptors (PRRs). These receptors help in labeling which pathogens are viruses, fungi, or bacteria. Moreover, the importance of siRNA and the innate immune system is to include more PRRs to help recognize different RNA structures. This makes it more likely for the siRNA to cause an immunostimulant response in the event of the pathogen.
Food
RNAi has been used to genetically engineer plants to produce lower levels of natural plant toxins. Such techniques take advantage of the stable and heritable RNAi phenotype in plant stocks. Cotton seeds are rich in dietary protein but naturally contain the toxic terpenoid product gossypol, making them unsuitable for human consumption. RNAi has been used to produce cotton stocks whose seeds contain reduced levels of delta-cadinene synthase, a key enzyme in gossypol production, without affecting the enzyme's production in other parts of the plant, where gossypol is itself important in preventing damage from plant pests.
Development efforts have successfully reduced the levels of allergens in tomato plants and fortification of plants such as tomatoes with dietary antioxidants. RNAi silencing of alpha-amylase have also been used to decrease Aspergillus flavus fungal growth in maize which would have otherwise contaminated the kernels with dangerous aflatoxins. Silencing lachrymatory factor synthase in onions have produced tearless onions and RNAi has been used in BP1 genes in rapeseeds to improve photosynthesis. SBEIIa and SBEIIb genes in wheat have been targeted in wheat in order to produce higher levels of amylose in order to improve bowel function, and Travella et al. 2006 employed RNAi for functional genomics an investigation of hexaploid bread races, while virus-induced gene silencing (VIGS, a subtype of RNAi) was used by Scofield et al. 2005 to investigate the mechanism of resistance provided by Lr21 against wheat leaf rust in hexaploid wheat.
Insecticides
RNAi is under development as an insecticide, employing multiple approaches, including genetic engineering and topical application. Cells in the midgut of some insects take up the dsRNA molecules in the process referred to as environmental RNAi. In some insects the effect is systemic as the signal spreads throughout the insect's body (referred to as systemic RNAi).
Animals exposed to RNAi at doses millions of times higher than anticipated human exposure levels show no adverse effects. RNAi has varying effects in different species of Lepidoptera (butterflies and moths).
Drosophila spp., Bombyx mori, Locusta spp., Spodoptera spp., Tribolium castaneum, Nilaparvata lugens, Helicoverpa armigera, and Apis mellifera are models that have been widely used to learn how RNAi works within particular taxa of insects. Musca domestica has two Ago2 genes and Glossina morsitans three, as found by Lewis et al. 2016 and Hain et al. 2010. In the case of the miRNA pathway Diuraphis noxia has two Ago1s, M. domestica two Dcr1s, Acyrthosiphon pisum two each of Ago1 and Loqs and Dcr1 and four of Pasha. While in the piRNA, G. morsitans and A. pisum have two or three Ago3s each. This has led to identification of future insecticide development targets, and the modes of action and reasons for insecticide resistance of other insecticides.
Transgenic plants
Transgenic crops have been made to express dsRNA, carefully chosen to silence crucial genes in target pests. These dsRNAs are designed to affect only insects that express specific gene sequences. As a proof of principle, in 2009 a study showed RNAs that could kill any one of four fruit fly species while not harming the other three.
Topical
Alternatively dsRNA can be supplied without genetic engineering. One approach is to add them to irrigation water. The molecules are absorbed into the plants' vascular system and poison insects feeding on them. Another approach involves spraying dsRNA like a conventional pesticide. This would allow faster adaptation to resistance. Such approaches would require low cost sources of dsRNAs that do not currently exist.
Functional genomics
Approaches to the design of genome-wide RNAi libraries can require more sophistication than the design of a single siRNA for a defined set of experimental conditions. Artificial neural networks are frequently used to design siRNA libraries and to predict their likely efficiency at gene knockdown. Mass genomic screening is widely seen as a promising method for genome annotation and has triggered the development of high-throughput screening methods based on microarrays.
Genome-scale screening
Genome-scale RNAi research relies on high-throughput screening (HTS) technology. RNAi HTS technology allows genome-wide loss-of-function screening and is broadly used in the identification of genes associated with specific phenotypes. This technology has been hailed as a potential second genomics wave, following the first genomics wave of gene expression microarray and single nucleotide polymorphism discovery platforms.
One major advantage of genome-scale RNAi screening is its ability to simultaneously interrogate thousands of genes. With the ability to generate a large amount of data per experiment, genome-scale RNAi screening has led to an explosion of data generation rates. Exploiting such large data sets is a fundamental challenge, requiring suitable statistics/bioinformatics methods. The basic process of cell-based RNAi screening includes the choice of an RNAi library, robust and stable cell types, transfection with RNAi agents, treatment/incubation, signal detection, analysis and identification of important genes or therapeutical targets.
History
RNAi discovery
The process of RNAi was referred to as "co-suppression" and "quelling" when observed prior to the knowledge of an RNA-related mechanism. The discovery of RNAi was preceded first by observations of transcriptional inhibition by antisense RNA expressed in transgenic plants, and more directly by reports of unexpected outcomes in experiments performed by plant scientists in the United States and the Netherlands in the early 1990s. In an attempt to alter flower colors in petunias, researchers introduced additional copies of a gene encoding chalcone synthase, a key enzyme for flower pigmentation into petunia plants of normally pink or violet flower color. The overexpressed gene was expected to result in darker flowers, but instead caused some flowers to have less visible purple pigment, sometimes in variegated patterns, indicating that the activity of chalcone synthase had been substantially decreased or became suppressed in a context-specific manner. This would later be explained as the result of the transgene being inserted adjacent to promoters in the opposite direction in various positions throughout the genomes of some transformants, thus leading to expression of antisense transcripts and gene silencing when these promoters are active. Another early observation of RNAi came from a study of the fungus Neurospora crassa, although it was not immediately recognized as related. Further investigation of the phenomenon in plants indicated that the downregulation was due to post-transcriptional inhibition of gene expression via an increased rate of mRNA degradation. This phenomenon was called co-suppression of gene expression, but the molecular mechanism remained unknown.
Not long after, plant virologists working on improving plant resistance to viral diseases observed a similar unexpected phenomenon. While it was known that plants expressing virus-specific proteins showed enhanced tolerance or resistance to viral infection, it was not expected that plants carrying only short, non-coding regions of viral RNA sequences would show similar levels of protection. Researchers believed that viral RNA produced by transgenes could also inhibit viral replication. The reverse experiment, in which short sequences of plant genes were introduced into viruses, showed that the targeted gene was suppressed in an infected plant. This phenomenon was labeled "virus-induced gene silencing" (VIGS), and the set of such phenomena were collectively called post transcriptional gene silencing.
After these initial observations in plants, laboratories searched for this phenomenon in other organisms. The first instance of RNA silencing in animals was documented in 1996, when Guo and Kemphues observed that, by introducing sense and antisense RNA to par-1 mRNA in Caenorhabditis elegans caused degradation of the par-1 message. It was thought that this degradation was triggered by single-stranded RNA (ssRNA), but two years later, in 1998, Fire and Mello discovered that this ability to silence the par-1 gene expression was actually triggered by double-stranded RNA (dsRNA). Craig C. Mello and Andrew Fire's 1998 Nature paper reported a potent gene silencing effect after injecting double stranded RNA into C. elegans. In investigating the regulation of muscle protein production, they observed that neither mRNA nor antisense RNA injections had an effect on protein production, but double-stranded RNA successfully silenced the targeted gene. As a result of this work, they coined the term RNAi. This discovery represented the first identification of the causative agent for the phenomenon. Fire and Mello were awarded the 2006 Nobel Prize in Physiology or Medicine.
RNAi therapeutics
Just after Fire and Mello's ground-breaking discovery, Elbashir et al. discovered, by using synthetically made small interfering RNA (siRNA), it was possible to target the silencing of specific sequences in a gene, rather than silencing the entire gene. Only a year later, McCaffrey and colleagues demonstrated that this sequence-specific silencing had therapeutic applications by targeting a sequence from the Hepatitis C virus in transgenic mice. Since then, multiple researchers have been attempting to expand the therapeutic applications of RNAi, specifically looking to target genes that cause various types of cancer. By 2006, the first applications to reach clinical trials were in the treatment of macular degeneration and respiratory syncytial virus. Four years later the first-in-human Phase I clinical trial was started, using a nanoparticle delivery system to target solid tumors.
The FDA approved the first siRNA-based drug (patisiran) in 2018. Givosiran and lumasiran later won FDA approval for the treatment of AHP and PH1 in 2019 and 2020, respectively. Inclisiran received EMA approval in 2020 for the treatment of high cholesterol and is currently under review by the FDA.
| Biology and health sciences | Molecular biology | Biology |
34740042 | https://en.wikipedia.org/wiki/Laminar%20flow%20reactor | Laminar flow reactor | A laminar flow reactor (LFR) is a type of chemical reactor that uses laminar flow to control reaction rate, and/or reaction distribution. LFR is generally a long tube with constant diameter that is kept at constant temperature. Reactants are injected at one end and products are collected and monitored at the other. Laminar flow reactors are often used to study an isolated elementary reaction or multi-step reaction mechanism.
Overview
Laminar flow reactors employ the characteristics of laminar flow to achieve various research purposes. For instance, LFRs can be used to study fluid dynamics in chemical reactions, or they can be utilized to generate special chemical structures such as carbon nanotubes. One feature of the LFR is that the residence time (The time interval during which the chemicals stay in the reactor) of the chemicals in the reactor can be varied by either changing the distance between the reactant input point and the point at which the product/sample is taken, or by adjusting the velocity of the gas/fluid. Therefore the benefit of a laminar flow reactor is that the different factors that may affect a reaction can be easily controlled and adjusted throughout an experiment.
Means of analyzing reactants in LFR
Means of analyzing the reaction include using a probe that enters into the reactor; or more accurately, sometimes one can utilize non-intrusive optical methods (e.g. use spectrometer to identify and analyze contents) to study reactions in the reactor. Moreover, taking the entire sample of the gas/fluid at the end of the reactor and collecting data may be useful as well. Using methods mentioned above, various data such as concentration, flow velocity etc. can be monitored and analyzed.
Flow velocity in LFR
Fluids or gases with controlled velocity pass through a laminar flow reactor in a fashion of laminar flow. That is, streams of fluids or gases slide over each other like cards. When analyzing fluids with the same viscosity ("thickness" or "stickiness") but different velocity, fluids are typically characterized into two types of flows: laminar flow and turbulent flow. Compared to turbulent flow, laminar flow tends to have a lower velocity and is generally at a lower Reynolds number. Turbulent flow, on the other hand, is irregular and travels at a higher speed. Therefore the flow velocity of a turbulent flow on one cross section is often assumed to be constant, or "flat". The "non-flat" flow velocity of laminar flow helps explain the mechanism of an LFR. For the fluid/gas moving in an LFR, the velocity near the center of the pipe is higher than the fluids near the wall of the pipe. Thus, the velocity distribution of the reactants tends to decrease from the center to the wall.
Residence time distribution (RTD)
The velocity near the center of the pipe is higher than the fluids near the wall of the pipe. Thus, the velocity distribution of the reactants tends to be higher in the center and lower on the side. Consider fluid being pumped through an LFR at constant velocity from the inlet, and the concentration of the fluid is monitored at the outlet. The graph of the residence time distribution should look like a negative slope with positive concavity. And the graph is modeled by the function: if is smaller than ; if is greater than or equal to . Notice that the graph has the value of zero initially, this is simply because it takes sometime for the substance to travel through the reactor. When the material is starting to reach the outlet, the concentration drastically increases, and it gradually decreases as time proceeds.
Characteristics
The laminar flows inside of a LFR has the unique characteristic of flowing in a parallel fashion without disturbing one another. The velocity of the fluid or gas will naturally decrease as it gets closer to the wall and farther from the center. Therefore the reactants have an increasing residence time in the LFR from the center to the side. A gradually increasing residence time gives researchers a clear layout of the reaction at different times.
Besides, when studying reactions in LFR, radial gradients in velocity, composition and temperature are significant. In other words, in other reactors where laminar flow is not significant, for instance, in a plug flow reactor, velocity of the object is assumed to be the same on one cross section since the flows are mostly turbulent. In a laminar flow reactor, velocity is significantly different at various points on the same cross section. Therefore the velocity differences throughout the reactor need to be taken into consideration when working with a LFR.
Research
Various researches pertaining to the modeling of LFR and formations of substances within a LFR have been done over the past decades. For instance, the formation of Single-walled carbon nanotube was investigated in a LFR. As another example, conversion from methane to higher hydrocarbons have been studied in a laminar flow reactor.
| Physical sciences | Chemical engineering | Chemistry |
27888690 | https://en.wikipedia.org/wiki/Erebidae | Erebidae | The Erebidae are a family of moths in the superfamily Noctuoidea. The family is among the largest families of moths by species count and contains a wide variety of well-known macromoth groups. The family includes the underwings (Catocala); litter moths (Herminiinae); tiger, lichen, footman and wasp moths (Arctiinae); tussock moths (Lymantriinae), including the arctic woolly bear moth (Gynaephora groenlandica); fruit-piercing moths (Calpinae and others); micronoctuoid moths (Micronoctuini); snout moths (Hypeninae); and zales, though many of these common names can also refer to moths outside the Erebidae (for example, crambid snout moths). Some of the erebid moths are called owlets.
The sizes of the adults range from among the largest of all moths (around wingspan in the white witch) to the smallest of the macromoths ( wingspan in some of the Micronoctuini). The coloration of the adults spans the full range of dull, drab, and camouflaged (e.g., Zale lunifera and litter moths) to vivid, contrasting, and colorful (e.g., Aganainae and tiger moths). The moths are found on all continents except Antarctica.
Subfamilies
Aganainae
Anobinae
Arctiinae – tiger, lichen, and wasp moths
Boletobiinae
Calpinae – piercing moths
Erebinae – underwings and kin
Eulepidotinae
Herminiinae – litter moths
Hypeninae – snout moths
Hypenodinae – includes the micronoctuoids
Hypocalinae
Lymantriinae – tussock moths
Pangraptinae
Rivulinae
Scolecocampinae
Scoliopteryginae – piercing moths
Tinoliinae
Toxocampinae
Description
Adult moths of Erebidae have quadrifid forewings and usually quadrifine hindwings, meaning that each wing includes a cubital vein that splits into four (explained further in the Classification section). The tribe Micronoctuini instead has bifine hindwings (meaning this same vein splits into two). Aside from this, adult Erebidae usually have an unscaled clypeofrons and broad forewings, and the hindwings often have patterns. Scoliopteryginae, Calpinae and some Erebinae have modified proboscises to pierce fruit skins. Lymantriinae and some Arctiina (Arctiinae) instead have a highly reduced proboscis. The overall appearance may be colourful (e.g. Arctiinae, Aganainae) or cryptic (e.g. Herminiinae).
Larvae are usually smooth in appearance, but larvae of Arctiinae (woolly bears) and Lymantriinae are hairy. Larvae of Arctiinae and Lymantriinae have fully developed prolegs, Aganainae and Herminiinae have fully developed or slightly reduced prolegs, and prolegs in some other subfamilies are reduced or absent as an adaptation to arboreal living (semi-loopers).
Ecology
Adults of some Erebidae pierce fruit to suck out juices (leading them to be called "fruit-piercing moths"), and those of Calyptra can also pierce mammalian skin to suck out blood (hence "vampire moths").
Larvae are mostly herbivorous, like most lepidopteran larvae, and different taxa prefer different plants. Lithosiini (Arctiinae) larvae are unusual in feeding on algae and lichens (hence "lichen moths"). Most Herminiinae larvae feed on dead or withered leaves (hence "litter moths"), rather than living leaves.
Classification
Among the Noctuoidea, the Erebidae can be broadly defined by the wing characteristics of the adults with support from phylogenetic studies. The cubital forewing vein, which runs outward from the base of a wing to the outer margin, splits into two (bifid), three (trifid), or four (quadrifid) veins from the medial area to the outer margin. These split veins are named M2, M3, CuA1, and CuA2 in order toward the inner margin. A trifid forewing has either a reduced or vestigial M2 vein or the M2 vein does not connect to the cubital veins, while M2 is as thick as M3 and connects or nearly connects to M3 in a quadrifid forewing. The same splitting of the hindwing cubital vein has analogous terms bifine, trifine, and quadrifine. The Erebidae typically have quadrifid forewings and quadrifine hindwings, though the Micronoctuini are exceptional with their bifine hindwings. Among the related families, most Erebidae are quadrifid moths like the Euteliidae, Nolidae, and Noctuidae and unlike the trifid Oenosandridae and Notodontidae. And among the quadrifid moths, the Erebidae have quadrifine hindwings like the typical Nolidae and Euteliidae and unlike the typical Noctuidae.
Phylogenetic studies in the present century have helped to clarify the relationships between the structurally diverse lineages within the Noctuoidea and within the Erebidae. Morphological studies had led to a classification in which the monophyletic Arctiinae, Lymantriinae, and Micronoctuini were treated as families, and the other erebid lineages were largely grouped within the Noctuidae. Recent studies combining genetic characteristics with the morphological ones revealed that the former Noctuidae were paraphyletic, and some of the lineages within the Noctuidae were more closely related to the Arctiinae, Lymantriinae, and Micronoctuini families than to the other lineages within the Noctuidae.
The determination of these phylogenetic relationships has led to the present classification scheme in which several clades were rearranged while kept mostly intact and others were split apart. The Erebidae are one monophyletic family among six in the Noctuoidea. A more strictly defined family Noctuidae is also monophyletic, but the family lacks the quadrifine moths now placed as part of the Erebidae. Some subfamilies of the Noctuidae, such as the Herminiinae, were moved as a whole to Erebidae. Other subfamilies, including the Acontiinae and Calpinae, were each split apart. The Arctiinae became an erebid subfamily placed next to the closely related Herminiinae. The Lymantriinae became another erebid subfamily placed near the Pangraptinae. The rank of the Micronoctuini was changed from family to tribe to include the clade as a lineage within the Hypenodinae. The Erebidae are currently divided into 18 subfamilies, some of which are strongly supported by phylogenetic analysis and may persist through further study, while others are weakly supported and may be redefined again.
| Biology and health sciences | Lepidoptera | Animals |
5301750 | https://en.wikipedia.org/wiki/Pteraspis | Pteraspis | Pteraspis (from 'wing' or 'fin' and 'shield') is an extinct genus of pteraspidid heterostracan jawless fish. It lived from the Lochkovian to Eifelian epochs of the Devonian period in what is now Brazil (Eifelian Maecuru Formation), Britain (Lochkovian Ditton Group), Ukraine (Lochkovian Ivane Suite, Pragian Babin Sandstone) and Belgium.
Description
Like other heterostracan fishes, Pteraspis had a protective armored plating covering the front of its body. Though lacking fins other than its lobed tail, it is thought to have been a good swimmer due to stiff, wing-like protrusions derived from the armoured plates over its gills. This, along with the horn-like rostrum, made Pteraspis very streamlined in shape, which is good for swimming. Pteraspis also had some stiff spikes on its back, possibly an additional form of protection against predators. It is thought to have fed from shoals of plankton just under the ocean surface. Some records are found in association with marine fossils, while some others are found in freshwater environment.
Pteraspis grew to an estimated length of .
Gallery
| Biology and health sciences | Prehistoric agnathae and early chordates | Animals |
5302698 | https://en.wikipedia.org/wiki/Brosimum%20alicastrum | Brosimum alicastrum | Brosimum alicastrum, commonly known as breadnut, Maya nut or ramon, and many others, is a tree species in the family Moraceae of flowering plants, whose other genera include figs and mulberries.
Two subspecies are commonly recognized:
B. a. alicastrum
B. a. bolivarense (Pittier) C.C.Berg
Description
Brosimum alicastrum can be monoecious, dioecious or hermaphroditic, changing from female to male as they age. Birds and bats are responsible for the dispersion of the seeds. A tree can produce of fruits per year. It stays productive for 120–150 years.
The tree can grow up to 45 m (150 ft) in height and up to in diameter.
It starts producing flowers and fruits when the tree's trunk reaches high.
When planted from seed in full sun, fruiting can start at 3.5 years.
Distribution and habitat
This tree is found on the west coast of central Mexico and in southern Mexico (Yucatán, Campeche), Guatemala, El Salvador, the Caribbean, and the Amazon basin. Large stands occur in moist lowland tropical forests at elevation (especially 125–800 m), in humid areas with annual rainfall of , and average temperatures of 24 °C (75 °F).
The Maya nut fruit disperses on the ground at different times throughout its range. It has a large seed covered by a thin, citrus-flavored, orange-colored skin favored by a number of forest creatures.
Cultivation
Breadnut may have formed a part of the diet of the pre-Columbian Maya of the lowlands region in Mesoamerica, although to what extent has been a matter of some debate among historians and archaeologists: no verified remains or illustrations of the fruit have been found at any Mayan archaeological sites.
It has been claimed in several publications by Dennis E. Puleston to have been a staple food in the Maya diet. Puleston demonstrated a strong correlation between ancient Maya settlement patterns and the distribution of relic stands of ramon trees.
Other research has downplayed the Maya nut's significance. In the modern era, it has been marginalized as a source of nutrition and has often been characterized as a famine food.
The tree lends its name to the Maya archaeological sites of Iximché and Topoxte, both in Guatemala and Tamuin (reflecting the Maya origin of the Huastec peoples). It is one of the 20 dominant species of the Maya forest. Of the dominant species, it is the only one that is wind-pollinated. It is also found in traditional Maya forest gardens.
A high density of seeds during the seedling offsets a reduced viability of the young plants and therefore enables a good yield. Seed storage is a common issue in seedling production. Long storage adversely affects the germination rate, for example after three weeks it decreases by 10%. Refrigeration is not a solution as it risks killing the seeds.
Uses
The Maya nut is high in fiber, calcium, potassium, iron, zinc, protein and B vitamins. It has a low glycemic index (<50) and is very high in antioxidants and prebiotic fiber. The fresh seeds can be cooked and eaten or can be set out to dry in the sun and roasted and milled into a chocolatey tasting powder. Stewed, the nut tastes like mashed potato; roasted, it tastes like chocolate or coffee. It can be prepared in numerous other dishes. In Petén, Guatemala, the breadnut is cultivated for exportation and local consumption as powder, for hot beverages, and bread.
The large seed is edible and can be boiled or dried and ground into a meal for porridge or flatbread.
Other uses
Breadnut leaves are commonly used as forage for livestock during the dry season in Central America. The fruits and seeds are also used to feed all kinds of animals.
Brosimum alicastrum can be used for carbon farming as a nut crop or fodder.
It is an oxalogene tree. It can therefore undertake a bacterial-fungal endosymbiosis which assists the oxalate-carbonate pathway (OCP) and especially the chemical reaction of biomineralization, and in this case biocalcification (to produce CaCO3 from CO2 and to store it in the soils). This tree would therefore act as a carbon sink, while providing resources for both humans and animals.
This was first shown by a biogeochemist Eric Verrechia, researcher at University of Lausanne in 2006.
The species can be used to restore damaged soils. It can prevent erosion and act as a wind barrier. The tree tolerates poor, damaged, dried or salty soils and it requires few inputs after its planting. Furthermore, its oxalogenic activity increases the pH and the amount of organic matter in the soil once well implemented in the agricultural system. This leads to an increased fertility thanks to a buffer effect. Some research projects are currently on-going to develop this crop in its current distribution area.
In culture
The name "breadnut" probably arose because the seeds can be ground to produce bread.
The plant is known by a range of names in indigenous Mesoamerican and other languages, including: ojoche, ojite, ojushte, ujushte, ujuxte, capomo, mojo, ox, iximche, masica in Honduras, uje in the state of Michoacan Mexico, mojote in Jalisco, in Haitian Creole and chataigne in Trinidadian Creole. In the Caribbean coast of Colombia it is called guaímaro or guaymaro.
| Biology and health sciences | Nuts | Plants |
5305778 | https://en.wikipedia.org/wiki/Monopotassium%20glutamate | Monopotassium glutamate | Monopotassium glutamate (MPG) is the compound with formula KC5H8NO4. It is a potassium salt of glutamic acid.
It has the E number E622 and is used in foods as a flavor enhancer. It is a non-sodium MSG alternative.
| Physical sciences | Glutamates | Chemistry |
31692117 | https://en.wikipedia.org/wiki/Raspberry%20Pi | Raspberry Pi | Raspberry Pi () is a series of small single-board computers (SBCs) developed in the United Kingdom. The original Raspberry Pi computer was developed by the Raspberry Pi Foundation in association with Broadcom. Since 2012, all Raspberry Pi products have been developed by Raspberry Pi Ltd, which began as a wholly-owned subsidiary of the Foundation.
The Raspberry Pi project originally leaned towards the promotion of teaching basic computer science in schools. The original model became more popular than anticipated, selling outside its target market for diverse uses such as robotics, home automation, industrial automation, and by computer and electronic hobbyists, because of its low cost, modularity, open design, and its adoption of the HDMI and USB standards.
The Raspberry Pi became the best-selling British computer in 2015, when it surpassed the ZX Spectrum in unit sales.
Origins and company history
The Raspberry Pi Foundation was created as a private company limited by guarantee in 2008, and was registered as a charity in 2009 by people at the University of Cambridge Computer Laboratory who had noticed a decline in the number and skills of young people applying for computer science courses.
In 2012, after the release of the second board type, the Raspberry Pi Foundation set up a new entity responsible for developing their computers, named Raspberry Pi (Trading) Ltd, and installed Eben Upton (one of the 2008 group) as CEO. The Foundation was rededicated as an educational charity for promoting the teaching of basic computer science in schools and developing countries.
In 2021, Raspberry Pi (Trading) Ltd changed its name to Raspberry Pi Ltd. Its newly-formed parent company, Raspberry Pi Holdings Ltd, became a public company in June 2024, launching on the London Stock Exchange where it trades with the stock symbol RPI.
Most Raspberry Pis are made in a Sony factory in Pencoed, Wales, while others are made in China and Japan.
Series and generations
There are three series of Raspberry Pi, and several generations of each have been released. Raspberry Pi SBCs feature a Broadcom system on a chip (SoC) with an integrated ARM-compatible central processing unit (CPU) and on-chip graphics processing unit (GPU), while Raspberry Pi Pico has a RP2040 system on chip with an integrated ARM-compatible central processing unit (CPU).
Raspberry Pi
The first-generation Raspberry Pi Model B was released in February 2012, followed by the simpler and cheaper Model A.
Raspberry Pi Model B+, an improved design, was released in 2014. These first-generation boards feature ARM11 processors, are approximately credit-card sized, and represent the standard mainline form factor. The A+ and an improved B model were released within a year. A "Compute Module" was released in April 2014 for embedded applications.
The Raspberry Pi 2 B was released in February 2015 and initially featured a 900 MHz 32-bit quad-core ARM Cortex-A7 processor with 1 GB RAM. Revision 1.2 features a 900 MHz 64-bit quad-core ARM Cortex-A53 processor (the same as that in the Raspberry Pi 3 Model B, but underclocked to 900 MHz).
The Raspberry Pi 3 Model B was released in February 2016 with a 1.2 GHz 64-bit quad core ARM Cortex-A53 processor, on-board 802.11n Wi-Fi, Bluetooth and USB boot capabilities.
The Raspberry Pi 3 Model B+ was launched on Pi Day 2018 with a faster 1.4 GHz processor, a three-times faster Gigabit Ethernet (throughput limited to ca. 300 Mbit/s by the internal USB 2.0 connection), and 2.4 / 5 GHz dual-band 802.11ac Wi-Fi (100 Mbit/s). Other features are Power over Ethernet (PoE) (with the add-on PoE HAT), USB boot and network boot (an SD card is no longer required).
The Raspberry Pi 3 Model A+ was launched in November 2018 as a similar board to the first Model A. It has a 1.4 GHz 64-bit quad-core processor, with 2.4 GHz dual-band and 5 GHz wireless LAN & Bluetooth 4.2. It also has a 40-pin GPIO header, 512 MB of DDR2 RAM, is powered by 5V of DC power via microUSB. A full-size HDMI port is used for connectivity, and one USB 2.0 port is on the board.
The Raspberry Pi 4 Model B was released in June 2019 with a 1.5 GHz 64-bit quad core ARM Cortex-A72 processor, on-board 802.11ac Wi-Fi, Bluetooth 5, full gigabit Ethernet (throughput not limited), two USB 2.0 ports, two USB 3.0 ports, 1, 2, 4, or 8 GB of RAM, and dual-monitor support via a pair of micro HDMI (HDMI Type D) ports for up to 4K resolution. The version with 1 GB RAM has been abandoned and the prices of the 2 GB version have been reduced. The 8 GB version has a revised circuit board. The Raspberry Pi 4 is also powered via a USB-C port, enabling additional power to be provided to downstream peripherals, when used with an appropriate PSU. But the Pi can only be operated with 5 volts and not 9 or 12 volts like other mini computers of this class. The initial Raspberry Pi 4 board had a design flaw where third-party e-marked USB cables, such as those used on MacBooks, incorrectly identify it and refuse to provide power. Tom's Hardware tested 14 different cables and found that 11 of them turned on and powered the Pi without issue. The design flaw was fixed in revision 1.2 of the board, released in late 2019. In mid-2021, Pi 4 B models appeared with the improved Broadcom BCM2711C0. The manufacturer is now using this chip for the Pi 4 B and Pi 400. However, the clock frequency of the Pi 4 B was not increased in the factory.
The Raspberry Pi 400 was released in November 2020. A modern example of a keyboard computer, it features 4 GB of LPDDR4 RAM on a custom board derived from the existing Raspberry Pi 4 combined with a keyboard in a single case. The case was derived from that of the Raspberry Pi Keyboard. A robust cooling solution (i.e. a broad metal plate) and an upgraded switched-mode power supply allow the Raspberry Pi 400's Broadcom BCM2711C0 processor to be clocked at 1.8 GHz, which is 20% faster than the Raspberry Pi 4 upon which it is based.
The Raspberry Pi 5 was announced in September 2023. It uses a 2.4 GHz quad-core 64-bit ARM Cortex-A76 CPU and a VideoCore VII GPU, with the improvements in hardware and software reportedly making the Pi 5 more than twice as powerful as the Pi 4. It has an I/O controller designed in-house, a power button, and an RTC chip (which requires an external battery). At launch, the Pi 5 was available with either 4 or 8 GB of RAM, at US$60 and US$80; a 2 GB variant was released in August 2024 at US$50. The Pi 5 lacks a 3.5 millimeter audio jack, so Bluetooth, HDMI, USB audio or an Audio HAT are the options for audio output.
In December 2024, Raspberry Pi introduced the keyboard-based Raspberry Pi 500, successor to the Pi 400.
In January 2025, the 16GB Raspberry Pi 5 was released, sold for $120.
Raspberry Pi Zero
The Raspberry Pi Zero with smaller size and reduced input/output (I/O) and general-purpose input/output (GPIO) capabilities was released in November 2015 for US$5.
The Raspberry Pi Zero v1.3 was released in May 2016, which added a camera connector.
The Raspberry Pi Zero W was launched in February 2017, a version of the Zero with Wi-Fi and Bluetooth capabilities, for US$10.
The Raspberry Pi Zero WH was launched in January 2018, a version of the Zero W with pre-soldered GPIO headers.
The Raspberry Pi Zero 2 W was launched in October 2021, a version of the Zero W with a system in a package (SiP) designed by Raspberry Pi and based on the Raspberry Pi 3. In contrast to the older Zero models, the Pi Zero 2 W is 64-bit capable. The price is around US$15.
Raspberry Pi Pico
Raspberry Pi Pico was released in January 2021 with a retail price of $4. It was Raspberry Pi's first board based upon a single microcontroller chip; the RP2040, which was designed by Raspberry Pi in the UK. The Pico has 264 KB of RAM and 2 MB of flash memory. It is programmable in C, C++, Assembly, MicroPython, CircuitPython and Rust. Raspberry Pi has partnered with Adafruit, Pimoroni, Arduino and SparkFun to build accessories for Raspberry Pi Pico and variety of other boards using RP2040 Silicon Platform. Rather than perform the role of general purpose computer (like the others in the range) it is designed for physical computing, similar in concept to an Arduino.
The Raspberry Pi Pico W was launched in June 2022, a version of the Pico with 802.11n Wi-Fi capability, for US$6. The CYW43439 wireless chip in the Pico W also supports Bluetooth, but the capability was not enabled at launch.
The Raspberry Pi Pico 2 was launched in August 2024 with a retail price of $5, based on a new RP2350 Arm/RISC-V microcontroller. The Pico 2 has 520 KB of RAM and 4 MB of flash memory and is hardware and software compatible with the original Pico.
The Raspberry Pi Pico 2 W was released in November 2024, with a retail price of $7 and using the RP2350 microcontroller. It has 4 MB of on-board flash memory to store code, while the RP2350 features 520 KB of on-chip SRAM. As for wireless capabilities, the Pico 2 W supports Wi-Fi (2.4GHz 802.11n) and Bluetooth 5.2.
Model comparison
Hardware
The Raspberry Pi hardware has evolved through several versions that feature variations in the type of the central processing unit, amount of memory capacity, networking support, and peripheral-device support.
This block diagram describes models B, B+, A and A+. The Pi Zero models are similar, but lack the Ethernet and USB hub components. The Ethernet adapter is internally connected to an additional USB port. In Model A, A+, and the Pi Zero, the USB port is connected directly to the system on a chip (SoC). On the Pi 1 Model B+ and later models the USB/Ethernet chip contains a five-port USB hub, of which four ports are available, while the Pi 1 Model B only provides two. On the Pi Zero, the USB port is also connected directly to the SoC, but it uses a micro USB (OTG) port. Unlike all other Pi models, the 40 pin GPIO connector is omitted on the Pi Zero, with solderable through-holes only in the pin locations. The Pi Zero WH remedies this.
Processor speed ranges from 700 MHz to 2.4 GHz for the Pi 5; on-board memory ranges from 256 MB to 8 GB random-access memory (RAM), with only the Raspberry Pi 4 and the Raspberry Pi 5 having more than 1 GB. Secure Digital (SD) cards in MicroSDHC form factor (SDHC on early models) are used to store the operating system and program memory, however some models also come with onboard eMMC storage and the Raspberry Pi 4 can also make use of USB-attached SSD storage for its operating system. The boards have one to five USB ports. For video output, HDMI and composite video are supported, with a standard 3.5 mm tip-ring-sleeve jack carrying mono audio together with composite video. Lower-level output is provided by a number of GPIO pins, which support common protocols like I²C. The B-models have an 8P8C Ethernet port and the Pi 3, Pi 4 and Pi Zero W have on-board Wi-Fi 802.11n and Bluetooth.
Processor
The Broadcom BCM2835 SoC used in the first generation Raspberry Pi includes a RISC-based 700 MHz 32-bit ARM1176JZF-S processor, VideoCore IV graphics processing unit (GPU), and RAM. It has a level 1 (L1) cache of 16 KB and a level 2 (L2) cache of 128 KB. The level 2 cache is used primarily by the GPU. The SoC is stacked underneath the RAM chip, so only its edge is visible. The ARM1176JZ(F)-S is the same CPU used in the original iPhone, although at a higher clock rate, and mated with a much faster GPU.
The earlier V1.1 model of the Raspberry Pi 2 used a Broadcom BCM2836 SoC with a 900 MHz 32-bit, quad-core ARM Cortex-A7 processor, with 256 KB shared L2 cache. The Raspberry Pi 2 V1.2 was upgraded to a Broadcom BCM2837 SoC with a 1.2 GHz 64-bit quad-core ARM Cortex-A53 processor, the same one which is used on the Raspberry Pi 3, but underclocked (by default) to the same 900 MHz CPU clock speed as the V1.1. The BCM2836 SoC is no longer in production as of late 2016.
The Raspberry Pi 3 Model B uses a Broadcom BCM2837 SoC with a 1.2 GHz 64-bit quad-core ARM Cortex-A53 processor, with 512 KB shared L2 cache. The Model A+ and B+ are 1.4 GHz
The Raspberry Pi 4 uses a Broadcom BCM2711 SoC with a 1.5 GHz (later models: 1.8 GHz) 64-bit quad-core ARM Cortex-A72 processor, with 1 MB shared L2 cache. Unlike previous models, which all used a custom interrupt controller poorly suited for virtualisation, the interrupt controller on this SoC is compatible with the ARM Generic Interrupt Controller (GIC) architecture 2.0, providing hardware support for interrupt distribution when using ARM virtualisation capabilities. The VideoCore IV of the previous models has also been replaced with a VideoCore VI running at 500 MHz.
The Raspberry Pi Zero and Zero W use the same Broadcom BCM2835 SoC as the first generation Raspberry Pi, although now running at 1 GHz CPU clock speed.
The Raspberry Pi Zero 2 W uses the RP3A0-AU, which is a System-in-Package (SiP) design. The package contains a Broadcom BCM2710A1 processor, which is a 64-bit quad-core ARM Cortex-A53 clocked at 1 GHz, along with 512 MB of LPDDR2 SDRAM layered above. The Raspberry Pi 3 also uses the BCM2710A1 in its Broadcom BCM2837 SoC, but clocked at a higher 1.2 GHz.
The Raspberry Pi Pico uses the RP2040, a microcontroller containing dual ARM Cortex-M0+ cores running at 133 MHz, 6 banks of SRAM totalling 264 KB, and programmable IO for peripherals.
The Raspberry Pi 5 uses the Broadcom BCM2712 SoC, which is a chip designed in collaboration with Raspberry Pi. The SoC features a quad-core ARM Cortex-A76 processor clocked at 2.4 GHz, alongside a VideoCore VII GPU clocked at 800 MHz. The BCM2712 SoC also features support for cryptographic extensions for the first time on a Raspberry Pi model. Alongside the new processor and graphics unit, the monolithic design of the earlier BCM2711 has been replaced with a CPU and chipset (southbridge) architecture, as the IO functionality has been moved to the Raspberry Pi 5's custom RP1 chip.
Performance
While operating at 700 MHz by default, the first generation Raspberry Pi provided a real-world performance roughly equivalent to 0.041 GFLOPS. On the CPU level the performance is similar to a 300 MHz Pentium II of 1997–99. The GPU provides 1 Gpixel/s or 1.5 Gtexel/s of graphics processing or 24 GFLOPS of general purpose computing performance. The graphical capabilities of the Raspberry Pi are roughly equivalent to the performance of the Xbox of 2001.
Raspberry Pi 2 V1.1 included a quad-core Cortex-A7 CPU running at 900 MHz and 1 GB RAM. It was described as 4–6 times more powerful than its predecessor. The GPU was identical to the original. In parallelised benchmarks, the Raspberry Pi 2 V1.1 could be up to 14 times faster than a Raspberry Pi 1 Model B+.
The Raspberry Pi 3, with a quad-core Cortex-A53 processor, is described as having ten times the performance of a Raspberry Pi 1. Benchmarks showed the Raspberry Pi 3 to be approximately 80% faster than the Raspberry Pi 2 in parallelised tasks.
The Raspberry Pi 4, with a quad-core Cortex-A72 processor, is described as having three times the performance of a Raspberry Pi 3.
Overclocking
Most Raspberry Pi systems-on-chip can be overclocked to various degrees utilising the built in config.txt file in the boot sector of the Raspberry Pi OS. Overclocking is generally safe and does not automatically void the warranty of the Raspberry Pi; however, setting the "force_turbo" option to 1 bypasses voltage and temperature limits and voids the users warranty. In Raspberry Pi OS the overclocking options on boot can also be made by a software command running "sudo raspi-config" on Raspberry Pi 1, 2, and original 3B without voiding the warranty. In those cases the Pi automatically shuts the overclocking down if the chip temperature reaches ; an appropriately sized heat sink is needed to protect the chip from thermal throttling.
Newer versions of the firmware contain the option to choose between five overclock ("turbo") presets that, when used, attempt to maximise the performance of the SoC without impairing the lifetime of the board. This is done by monitoring the core temperature of the chip and the CPU load, and dynamically adjusting clock speeds and the core voltage. When the demand is low on the CPU or it is running too hot, the performance is throttled, but if the CPU has much to do and the chip's temperature is acceptable, performance is temporarily increased with CPU clock speeds of up to 1.1 GHz, depending on the board version and on which of the turbo settings is used.
The overclocking modes are:
In the highest (turbo) mode the SDRAM clock speed was originally 500 MHz, but this was later changed to 600 MHz because of occasional SD card corruption. Simultaneously, in high mode the core clock speed was lowered from 450 to 250 MHz, and in medium mode from 333 to 250 MHz.
The CPU of the first and second generation Raspberry Pi board did not require cooling with a heat sink or fan, even when overclocked, but the Raspberry Pi 3 may generate more heat when overclocked.
RAM
The early designs of the Raspberry Pi Model A and B boards included 256 MB of random-access memory (RAM). Of this, the early beta Model B boards allocated 128 MB to the GPU by default, leaving only 128 MB for the CPU. On the early 256 MB releases of models A and B, three different splits were possible. The default split was 192 MB for the CPU, which should be sufficient for standalone 1080p video decoding, or for simple 3D processing. 224 MB was for Linux processing only, with only a 1080p framebuffer, and was likely to fail for any video or 3D. 128 MB was for heavy 3D processing, possibly also with video decoding. In comparison, the Nokia 701 uses 128 MB for the Broadcom VideoCore IV.
The later Model B with 512 MB RAM, was released on 15 October 2012 and was initially released with new standard memory split files (arm256_start.elf, arm384_start.elf, arm496_start.elf) with 256 MB, 384 MB, and 496 MB CPU RAM, and with 256 MB, 128 MB, and 16 MB video RAM, respectively. But about one week later, the foundation released a new version of start.elf that could read a new entry in config.txt (gpu_mem=xx) and could dynamically assign an amount of RAM (from 16 to 256 MB in 8 MB steps) to the GPU, obsoleting the older method of splitting memory, and a single start.elf worked the same for 256 MB and 512 MB Raspberry Pis.
The Raspberry Pi 2 has 1 GB of RAM.
The Raspberry Pi 3 has 1 GB of RAM in the B and B+ models, and 512 MB of RAM in the A+ model. The Raspberry Pi Zero and Zero W have 512 MB of RAM.
The Raspberry Pi 4 is available with 1, 2, 4 or 8 GB of RAM. A 1 GB model was originally available at launch in June 2019 but was discontinued in March 2020, and the 8 GB model was introduced in May 2020. The 1 GB model returned in October 2021.
The Raspberry Pi 5 is available with 2, 4, 8 or 16 GB of RAM.
Networking
The Model A, A+ and Pi Zero have no Ethernet circuitry and are commonly connected to a network using an external user-supplied USB Ethernet or Wi-Fi adapter. On the the Ethernet port is provided by a built-in USB Ethernet adapter using the SMSC LAN9514 chip. The Raspberry Pi 3 and Pi Zero W (wireless) are equipped with 2.4 GHz WiFi 802.11n and Bluetooth 4.1 based on the Broadcom BCM43438 FullMAC chip with no official support for monitor mode (though it was implemented through unofficial firmware patching) and the Pi 3 also has a 10/100 Mbit/s Ethernet port. The Raspberry Pi 3B+ features dual-band IEEE 802.11b/g/n/ac WiFi, Bluetooth 4.2, and Gigabit Ethernet (limited to approximately 300 Mbit/s by the USB 2.0 bus between it and the SoC). The Raspberry Pi 4 has full gigabit Ethernet (throughput is not limited as it is not funnelled via the USB chip.)
Special-purpose features
The RPi Zero, RPi1A, RPi3A+ and RPi4 can be used as a USB device or "USB gadget", plugged into another computer via a USB port on another machine. It can be configured in multiple ways, such as functioning as a serial or Ethernet device. Although originally requiring software patches, this was added into the mainline Raspbian distribution in May 2016.
Raspberry Pi models with a newer chipset can boot from USB mass storage, such as from a flash drive. Booting from USB mass storage is not available in the original Raspberry Pi models, the Raspberry Pi Zero, the Raspberry Pi Pico, the Raspberry Pi 2 A models, and the Raspberry Pi 2 B models with versions lower than 1.2.
Peripherals
Although often pre-configured to operate as a headless computer, the Raspberry Pi may also optionally be operated with any generic USB computer keyboard and mouse. It may also be used with USB storage, USB to MIDI converters, and virtually any other device/component with USB capabilities, depending on the installed device drivers in the underlying operating system (many of which are included by default).
Other peripherals can be attached through the various pins and connectors on the surface of the Raspberry Pi.
Video
The video controller can generate standard modern TV resolutions, such as HD and Full HD, and higher or lower monitor resolutions as well as older NTSC or PAL standard CRT TV resolutions. As shipped (i.e., without custom overclocking) it can support the following resolutions: 640×350 EGA; 640×480 VGA; 800×600 SVGA; 1024×768 XGA; 1280×720 720p HDTV; 1280×768 WXGA variant; 1280×800 WXGA variant; 1280×1024 SXGA; 1366×768 WXGA variant; 1400×1050 SXGA+; 1600×1200 UXGA; 1680×1050 WXGA+; 1920×1080 1080p HDTV; 1920×1200 WUXGA.
Higher resolutions, up to 2048×1152, may work or even 3840×2160 at 15 Hz (too low a frame rate for convincing video). Allowing the highest resolutions does not imply that the GPU can decode video formats at these resolutions; in fact, the Raspberry Pis are known to not work reliably for H.265 (at those high resolutions), commonly used for very high resolutions (however, most common formats up to Full HD do work).
Although the Raspberry Pi 3 does not have H.265 decoding hardware, the CPU is more powerful than its predecessors, potentially fast enough to allow the decoding of H.265-encoded videos in software. The GPU in the Raspberry Pi 3 runs at higher clock frequencies of 300 MHz or 400 MHz, compared to previous versions which ran at 250 MHz.
The Raspberry Pis can also generate 576i and 480i composite video signals, as used on old-style (CRT) TV screens and less-expensive monitors through standard connectorseither RCA or 3.5 mm phono connector depending on model. The television signal standards supported are PAL-B/G/H/I/D, PAL-M, PAL-N, NTSC and NTSC-J.
Real-time clock
When booting, the time defaults to being set over the network using the Network Time Protocol (NTP). The source of time information can be another computer on the local network that does have a real-time clock, or to a NTP server on the internet. If no network connection is available, the time may be set manually or configured to assume that no time passed during the shutdown. In the latter case, the time is monotonic (files saved later in time always have later timestamps) but may be considerably earlier than the actual time. For systems that require a built-in real-time clock, a number of small, low-cost add-on boards with real-time clocks are available. The Raspberry Pi 5 is the first to include a real-time clock. If an external battery is not plugged in, the Pi 5 will use the Network Time Protocol, or will need to be set manually, as was the case in previous models.
The RP2040 microcontroller has a built-in real-time clock, but it can not be set without some form of user entry or network facility being added.
Connectors
Pi Pico
Pi Compute Module
Pi Zero
Model A
Model B
J8 header and general purpose input-output (GPIO)
Raspberry Pi 1 Models A+ and B+, Pi 2 Model B, Pi 3 Models A+, B and B+, Pi 4, and Pi Zero, Zero W, Zero WH and Zero W 2 have the same 40-pin pinout (designated J8 across all models). Raspberry Pi 1 Models A and B have only the first 26 pins. The J8 header is commonly referred to as the GPIO connector as a whole, even though only a subset of the pins are GPIO pins. In the Pi Zero and Zero W, the 40 GPIO pins are unpopulated, having the through-holes exposed for soldering instead. The Zero WH (Wireless + Header) has the header pins preinstalled.
Model B rev. 2 also has a pad (called P5 on the board and P6 on the schematics) of 8 pins offering access to an additional 4 GPIO connections. These GPIO pins were freed when the four board version identification links present in revision 1.0 were removed.
Models A and B provide GPIO access to the ACT status LED using GPIO 16. Models A+ and B+ provide GPIO access to the ACT status LED using GPIO 47, and the power status LED using GPIO 35.
Specifications
Simplified Model B changelog
Software
Operating systems
Raspberry Pi provides Raspberry Pi OS (formerly called Raspbian), a Debian-based Linux distribution for download, as well as third-party Ubuntu, Windows 10 IoT Core, RISC OS, LibreELEC (specialised media centre distribution) and specialised distributions for the Kodi media centre and classroom management. It promotes Python and Scratch as the main programming languages, with support for many other languages. The default firmware is closed source, while unofficial open source firmware is available. Many other operating systems can also run on the Raspberry Pi. The formally verified microkernel seL4 is also supported. There are several ways of installing multiple operating systems on one mSD card.
Other operating systems (not Linux- nor BSD-based)
Broadcom VCOS – Proprietary operating system which includes an abstraction layer designed to integrate with existing kernels, such as ThreadX (which is used on the VideoCore4 processor), providing drivers and middleware for application development. In the case of the Raspberry Pi, this includes an application to start the ARM processor(s) and provide the publicly documented API over a mailbox interface, serving as its firmware. An incomplete source of a Linux port of VCOS is available as part of the reference graphics driver published by Broadcom.
Haiku – an open source BeOS clone that has been compiled for the Raspberry Pi and several other ARM boards. Work on Pi 1 began in 2011, but only the Pi 2 will be supported.
HelenOS – a portable microkernel-based multiserver operating system; has basic Raspberry Pi support since version 0.6.0
Plan 9 from Bell Labs and Inferno (in beta)
QNX
RISC OS Pi (a cut-down version of RISC OS Pico, for 16 MB cards and larger for all models of Pi 1 & 2, has also been made available)
Ultibo Core – OS-less unikernel Run Time Library based on Free Pascal. Lazarus IDE (Windows with 3rd party ports to Linux and MacOS). Most Pi models supported.
Windows 10 IoT Core – a zero-price edition of Windows 10 offered by Microsoft that runs natively on the Raspberry Pi 2.
Other operating systems (Linux-based)
Alpine Linux – a Linux distribution based on musl and BusyBox, "designed for power users who appreciate security, simplicity and resource efficiency".
Android is available for non-commercial use from KonstaKANG
Arch Linux ARM – a port of Arch Linux for ARM processors; the Arch-based Manjaro is also available for ARM
arkOS – designed for website and email self-hosting
CentOS for Raspberry Pi 2 and later
Devuan
emteria.OS – an embedded, managed version of the Android operating system for professional fleet management
Fedora (supports Pi 2 and later since Fedora 25, Pi 1 is supported by some unofficial derivatives) and RedSleeve (a RHEL port) for Raspberry Pi 1
Gentoo Linux
Kali Linux – a Debian-derived distribution designed for digital forensics and penetration testing
MX Linux – based on Debian Stable and including antiX components, this OS is available in Xfce, from which KDE and Fluxbox versions can be produced
openSUSE, SUSE Linux Enterprise Server 12 SP2 and Server 12 SP3 (commercial support)
OpenWrt – a highly extensible Linux distribution for embedded devices (typically wireless routers). It supports Pi 1, 2, 3, 4 and Zero W.
Pop!_PI for Raspberry Pi 4 is a distribution of Pop!_OS 22.04
postmarketOS – distribution based on Alpine Linux, primarily developed for smartphones
RetroPie – an offshoot of Raspbian OS that uses Emulation Station as its frontend for RetroArch and other emulators like Mupen64 for retro gaming. Hardware like Freeplay tech can help replace Game boy internals with RetroPie emulation.
NixOS – a Linux distribution based on the purely functional package management system Nix. NixOS is composed using modules and packages defined in the Nixpkgs project.
Rocky Linux
Sailfish OS with Raspberry Pi 2 (due to use ARM Cortex-A7 CPU; Raspberry Pi 1 uses different ARMv6 architecture and Sailfish requires ARMv7.)
Slackware ARM – version 13.37 and later runs on the Raspberry Pi without modification. The 128–496 MB of available memory on the Raspberry Pi is at least twice the minimum requirement of 64 MB needed to run Slackware Linux on an ARM or i386 system. (Whereas the majority of Linux systems boot into a graphical user interface, Slackware's default user environment is the textual shell / command line interface.) The Fluxbox window manager running under the X Window System requires an additional 48 MB of RAM.
SolydXK – a light Debian-derived distro with Xfce
Tiny Core Linux – a minimal Linux operating system focused on providing a base system using BusyBox and FLTK. Designed to run primarily in RAM.
Tizen – a Linux-based mobile operating system that was backed by the Linux Foundation and was mainly developed and primarily used by Samsung
Trisquel, a fully free GNU/Linux distribution
Linux Q83 is a fast, secure, and innovative Linux distribution
Ubuntu-based: Lubuntu and Xubuntu
Void Linux – a rolling release Linux distribution which was designed and implemented from scratch, provides images based on musl or glibc
webOS Open Source Edition – an open source version of webOS
Other operating systems (BSD-based)
FreeBSD
NetBSD
OpenBSD (only on 64-bit platforms, such as Raspberry Pi 3/4)
Driver APIs
Raspberry Pi can use a VideoCore IV GPU via a binary blob, which is loaded into the GPU at boot time from the SD-card, and additional software, that initially was closed source. This part of the driver code was later released. However, much of the actual driver work is done using the closed source GPU code. Application software makes calls to closed source run-time libraries (OpenMAX IL, OpenGL ES or OpenVG), which in turn call an open source driver inside the Linux kernel, which then calls the closed source VideoCore IV GPU driver code. The API of the kernel driver is specific for these closed libraries. Video applications use OpenMAX IL, use OpenGL ES and use OpenVG, which both in turn use EGL. OpenMAX IL and EGL use the open source kernel driver in turn.
Vulkan driver
Raspberry Pi first announced it was working on a Vulkan driver in February 2020. A working Vulkan driver running Quake 3 at 100 frames per second on a 3B+ was revealed by a graphics engineer who had been working on it as a hobby project on 20 June. On 24 November 2020 Raspberry Pi announced that their driver for the Raspberry Pi 4 is Vulkan 1.0 conformant. Raspberry Pi Trading announced further driver conformance for Vulkan 1.1 and 1.2 on 26 October 2021 and 1 August 2022.
Firmware
The official firmware is a freely redistributable binary blob, that is proprietary software. A minimal proof-of-concept open source firmware is also available, mainly aimed at initialising and starting the ARM cores as well as performing minimal startup that is required on the ARM side. It is also capable of booting a very minimal Linux kernel, with patches to remove the dependency on the mailbox interface being responsive. It is known to work on Raspberry Pi 1, 2 and 3, as well as some variants of Raspberry Pi Zero.
Third-party application software
AstroPrint – AstroPrint's wireless 3D printing software can be run on the Pi 2.
C/C++ Interpreter Ch – Released 3 January 2017, C/C++ interpreter Ch and Embedded Ch are released free for non-commercial use for Raspberry Pi, ChIDE is also included for the beginners to learn C/C++.
Minecraft (Pi edition) – Released 11 February 2013 and support ended on 24 January 2016, a modified version that allows players to directly alter the world with computer code.
RealVNC – Since 28 September 2016, Raspbian includes RealVNC's remote access server and viewer software. This includes a new capture technology which allows directly rendered content (e.g. Minecraft, camera preview and omxplayer) as well as non-X11 applications to be viewed and controlled remotely.
Steam Link – On 13 December 2018, Valve released official Steam Link game streaming client for the Raspberry Pi 3 and 3 B+.
UserGate Web Filter – On 20 September 2013, Florida-based security vendor Entensys announced porting UserGate Web Filter to Raspberry Pi platform.
Software development tools
Algoid – for teaching programming to children and beginners.
Arduino IDE – for programming an Arduino.
BlueJ – for teaching Java to beginners.
C-STEM Studio – a platform for hands-on integrated learning of computing, science, technology, engineering, and mathematics (C-STEM) with robotics.
CircuitPython - an educational fork of MicroPython for microcontrollers and single-board computers
Erlang – a functional language for building concurrent systems with light-weight processes and message passing.
Greenfoot – Greenfoot teaches object orientation with Java. Create 'actors' which live in 'worlds' to build games, simulations, and other graphical programs.
Julia – an interactive and cross-platform programming language/environment, that runs on the Pi 1 and later. IDEs for Julia, such as Visual Studio Code, are available. | Technology | Specific hardware | null |
31694713 | https://en.wikipedia.org/wiki/Port%20of%20Jebel%20Ali | Port of Jebel Ali | Port of Jebel Ali, also known as Mina Jebel Ali, is a deep port located in Jebel Ali, Dubai, United Arab Emirates. Jebel Ali is the world's ninth busiest port, the largest man-made harbour, and the biggest and by far the busiest port in the Middle-East. Port Jebel Ali was constructed in the late 1970s to supplement the facilities at Port Rashid.
Geography
Jebel Ali port is located 35 km southwest of Dubai, in the Persian Gulf. The port is part of the Maritime Silk Road that runs from the Chinese coast to the south via the southern tip of India to Mombasa, from there through the Red Sea via the Suez Canal to the Mediterranean, there to the Upper Adriatic region to the northern Italian hub of Trieste with its rail connections to Central Europe, Eastern Europe and the North Sea.
History
Jebel Ali Port, credited to the efforts of Rashid bin Saeed Al-Maktoum, was constructed in the late 1970s and opened in 1979 to supplement the facilities at Port Rashid. It was inaugurated by Queen Elizabeth II on 26 February 1979. The village of Jebel Ali was constructed for port workers, and it has a population of 300 people. Covering over . It is home to over 5,000 companies from 120 countries of the world. With 67 berths and a size of , Jebel Ali is the world's largest man-made harbour and the biggest port in the Middle-East.
The port of Jebel Ali has become the port most frequently visited by ships of the United States Navy outside the United States. Due to the depth of the harbour and size of the port facilities, a and several ships of the accompanying battle group can be accommodated pier-side. Due to the frequency of these port visits, semi-permanent liberty facilities (referred to by service personnel as "The Sandbox") have been erected adjacent to the carrier berth.
Operations
Port Jebel Ali encompass over one million square metres of container yard. It also contains space for medium- and long-term general cargo storage, including seven Dutch barns with a total of almost 19 thousand square metres and 12 covered sheds covering with 90.5 square metres. In addition, Port Jebel Ali also consist of 960 thousand square metres of open storage.
Port Jebel Ali is linked to Dubai's expressway system and to the Dubai International Airport Cargo Village. The Cargo Village facilities capable of handling cargoes, making four-hour transit from ship to aircraft possible.
The DPA's commercial trucking service transport container and general cargo transport between Port Jebel Ali, Port Rashid, and the rest of UAE every day.
Jebel Ali port is one of DP World's flagship facilities and have been ranked as 9th in Top Container Port Worldwide having handled 7.62 million TEUs in 2005, which represents a 19% increase in throughput, over 2004. Jebel Ali Port was ranked 7th in the world's largest ports in 2007.
Jebel Ali port is managed by state-owned Dubai Ports World.
Expansion
The expansion of Jebel Ali port commenced in 2001, which is the master plan of the port. The project comprises 15 stages, which will be completed over the decade. Stage one was completed in 2007, which has increased the storage and handling capacity by 2.2 million TEUs and a Quay length of 1,200 m.
The entire project includes 2.4 km of new berths, the container yard behind the berths and the supporting infrastructure and buildings necessary for a fully functioning terminal. The new port will be on reclaimed land extending seaward from the existing port and situated to the west of the Jumeirah Palm Island complex.
The current plan is expected to multiply the total capacity of Jebel Ali port by more than seven, making it the world's biggest container port, surpassing the ports of Shanghai and Singapore.
Awards
On 9 April 2011 Port of Jebel Ali won the Golden Award for Best Seaport Overall from the Higher Committee for UAE Civil Seaports and Airports Security.
Incidents
In 2017, a big port crane started to collapse. It hit a vehicle and building as it collapsed. People were able to escape.
Gallery
| Technology | Specific piers and ports | null |
3935892 | https://en.wikipedia.org/wiki/Bay | Bay | A bay is a recessed, coastal body of water that directly connects to a larger main body of water, such as an ocean, a lake, or another bay. A large bay is usually called a gulf, sea, sound, or bight. A cove is a small, circular bay with a narrow entrance. A fjord is an elongated bay formed by glacial action.
The term embayment is also used for , such as extinct bays or freshwater environments.
A bay can be the estuary of a river, such as the Chesapeake Bay, an estuary of the Susquehanna River. Bays may also be nested within each other; for example, James Bay is an arm of Hudson Bay in northeastern Canada. Some large bays, such as the Bay of Bengal and Hudson Bay, have varied marine geology.
The land surrounding a bay often reduces the strength of winds and blocks waves. Bays may have as wide a variety of shoreline characteristics as other shorelines. In some cases, bays have beaches, which "are usually characterized by a steep upper foreshore with a broad, flat fronting terrace". Bays were significant in the history of human settlement because they provided easy access to marine resources like fisheries. Later they were important in the development of sea trade as the safe anchorage they provide encouraged their selection as ports.
Definition
The United Nations Convention on the Law of the Sea defines a bay as a well-marked indentation in the coastline, whose penetration is in such proportion to the width of its mouth as to contain land-locked waters and constitute more than a mere curvature of the coast. An indentation, however, shall not be regarded as a bay unless its area is as large as (or larger than) that of the semi-circle whose diameter is a line drawn across the mouth of that indentation — otherwise it would be referred to as a bight.
Types
Open bay — a bay that is widest at the mouth, flanked by headlands.
Enclosed bay — a bay whose mouth is narrower than its widest part, flanked by at least one peninsula.
Semi-enclosed bay — an open bay whose exit is made into narrower channels by one or more islands within its mouth.
Back-barrier bay — a semi-enclosed bay separated from open water by one or more barrier islands or spits.
Formation
There are various ways in which bays can form. The largest bays have developed through plate tectonics. As the super-continent Pangaea broke up along curved and indented fault lines, the continents moved apart and left large bays; these include the Gulf of Guinea, the Gulf of Mexico, and the Bay of Bengal, which is the world's largest bay.
Bays also form through coastal erosion by rivers and glaciers. A bay formed by a glacier is a fjord. Rias are created by rivers and are characterised by more gradual slopes. Deposits of softer rocks erode more rapidly, forming bays, while harder rocks erode less quickly, leaving headlands.
| Physical sciences | Oceanic and coastal landforms | null |
3937442 | https://en.wikipedia.org/wiki/Ahaetulla | Ahaetulla | Ahaetulla, commonly referred to as Asian vine snakes or Asian whip snakes, is a genus of colubrid snakes distributed throughout tropical Asia. They are considered by some scientists to be mildly venomous and are what is commonly termed as 'rear-fanged' or more appropriately, opisthoglyphous, meaning their enlarged teeth or fangs, intended to aid in venom delivery, are located in the back of the upper jaw, instead of in the front as they are in vipers or cobras. As colubrids, Ahaetulla do not possess a true venom gland or a sophisticated venom delivery system. The Duvernoy's gland of this genus, homologous to the venom gland of true venomous snakes, produces a secretion which, though not well studied, is considered not to be medically significant to humans.
Green-colored members of this genus are often referred to as green vine snakes. They are not to be confused with the "green vine snake" Oxybelis fulgidus, which convergently appears very similar but is found in Central and South America.
Etymology
The genus name Ahaetulla comes from the Sri Lankan Sinhalese words ahaetulla/ahata gulla/as gulla, meaning “eye plucker” or “eye picker”, because of the belief that they pluck out the eyes of humans, as first reported by the Portuguese traveler João Ribeiro in 1685.
Vernacular names
The Sinhala name "" or "eye-plucker" forms the taxonomic genus name. In Tamil, it is known as , in Malayalam it is known as , in Telugu it is known as , in Marathi, it is known as , and in Kannada it is known as . There are dozens of other vernacular names for this snake genus within its range.
Sinhala: ඇහැටුල්ලා (Pronounced: Aheatulla)
Telugu: పచ్చారిపాము; పసరికపాము
Bengali: লাউডগা.
Odia: ଲାଉଡଙ୍କିଆ
Kannada: ಹಸಿರು ಹಾವು, ಹಸಿರು ಬಳ್ಳಿ ಹಾವು.
Gujarati: લીલવણ, માળણ.
Marathi हरणटोळ, शेलाटी
Tamil: பச்சை பாம்பு
Malayalam: പച്ചില പാമ്പ്,കൺകൊത്തി
Taxonomy
Their closest relative is the monotypic genus Proahaetulla, which Ahaetulla diverged from an estimated 26.57 million years ago, during the mid-Oligocene. From here, the clade containing Proahaetulla and Ahaetulla is a sister group to the genus Dryophiops, and the clade containing all three of these genera is a sister group to the clade containing the bronzeback snakes (Dendrelaphis) and flying snakes (Chrysopelea). These relationships can be shown in the cladogram below, with possible paraphyletic Ahaetulla species noted:
In 2020, an analysis of Ahaetulla nasuta, Ahaetulla dispar, and Ahaetulla pulverulenta throughout their range found them to represent species complexes containing several undescribed or formerly synonymized species, leading to the description of A. borealis, A. farnsworthi, A. malabarica, A. travancorica, and A. sahyadrensis, as well as the resurrection of A. oxyrhyncha and A. isabellina. Ahaetulla nasuta and Ahaetulla pulverulenta, formerly considered to have much wider ranges, are now considered endemic to Sri Lanka.
Species
The taxonomy of vine snakes is not well-documented, and literature varies widely, but there are 20 currently accepted species in the genus Ahaetulla:
Ahaetulla anomala (Annandale, 1906) - Variable colored vine snake (possibly conspecific with A. oxyrhyncha)
Ahaetulla borealis Mallik, Srikanthan, Pal, Princia D'Souza, Shanker & Ganesh, 2020 - Northern Western Ghats vine snake
Ahaetulla dispar (Günther, 1864) - Günther's vine snake
Ahaetulla farnsworthi Mallik, Srikanthan, Pal, Princia D'Souza, Shanker & Ganesh, 2020 - Farnsworth's vine snake
Ahaetulla fasciolata (Fischer, 1885) - Speckle-headed whipsnake
Ahaetulla flavescens (Wall, 1910 - yellow whipsnake
Ahaetulla fronticincta (Günther, 1858) - Burmese vine snake
Ahaetulla fusca (Duméril, Bibron, & Duméril,1854) - dark whipsnake
Ahaetulla isabellina (Wall, 1910) - Wall's vine snake
Ahaetulla laudankia Deepak, Narayanan, Sarkar, Dutta & Mohapatra, 2019 - Laudankia vine snake
Mirza, Pattekar, Verma, Stuart, Purkayastha, Mohapatra & Patel, 2024 - long-snouted arboreal snake
Ahaetulla malabarica Mallik, Srikanthan, Pal, Princia D'Souza, Shanker & Ganesh, 2020 - Malabar vine snake
Ahaetulla mycterizans (Linnaeus, 1758) - Malayan green whipsnake
Ahaetulla nasuta (Lacépède, 1789) - Sri Lankan green vine snake
Ahaetulla oxyrhyncha (Bell, 1825) - Indian vine snake
Ahaetulla perroteti (Duméril & Bibron, 1854) - Nilgiri vine snake
Ahaetulla prasina (Boie, 1827) - Oriental whipsnake or Asian vine snake
Ahaetulla prasina prasina (Boie, 1827)
Ahaetulla prasina medioxima Lazell, 2002
Ahaetulla prasina preocularis (Taylor, 1922)
Ahaetulla prasina suluensis Gaulke, 1994
Ahaetulla pulverulenta (Duméril & Bibron, 1854) - Brown-speckled whipsnake
Ahaetulla rufusoculara Lam, Thu, Nguyen, Murphy, & Nguyen, 2021
Ahaetulla sahyadrensis Mallik, Srikanthan, Pal, Princia D'Souza, Shanker & Ganesh, 2020
Ahaetulla travancorica Mallik, Srikanthan, Pal, Princia D'Souza, Shanker & Ganesh, 2020 - Travancore vine snake
Several undescribed species (including the Southeast Asian Ahaetulla formerly assigned to A. nasuta) still likely remain in these complexes.
Geographic range
They are found from Sri Lanka and India to China and much of Southeast Asia. Sri Lanka and the Western Ghats of India are major hotspots of diversity for the genus, with at least 10 of the currently-described species being endemic to these regions.
Description
All Ahaetulla species are characterized by thin, elongated bodies, with extremely long tails and a sharply triangular shaped head. They are primarily green in color, but can vary quite a bit to yellows, oranges, greys, and browns. They can have black and/or white patterning, or can be solid in color. Their eyes are almost unique in the reptile world, having keen binocular vision and keyhole shaped pupils, being similar in this aspect with twig snakes, who also have keyhole shaped pupils.
Behavior
They are primarily diurnal and arboreal, living in humid rainforests. Their diet consists mainly of lizards, but sometimes frogs and rodents are also consumed. Ahaetulla fronticincta, however, feeds exclusively on fish, striking its prey from branches overhanging water. Ahaetulla venom is not considered to be dangerous to humans, but serves to cause paralysis in their fast moving prey choices. They are ovoviviparous.
In captivity
Ahaetulla species are not yet frequently captive bred, as are many of the more popular snakes in the reptile keeping hobby. They are suitable for more advanced keepers, requiring a humid arboreal habitat and a diet of lizards as they rarely switch to rodents. Without proper husbandry, they are prone to health issues and stress.
| Biology and health sciences | Snakes | Animals |
3938382 | https://en.wikipedia.org/wiki/Urine | Urine | Urine is a liquid by-product of metabolism in humans and in many other animals. In placental mammals, urine flows from the kidneys through the ureters to the urinary bladder and exits the urethra through the penis or vulva during urination. In other vertebrates, urine is excreted through the cloaca.
Urine contains water-soluble by-products of cellular metabolism that are rich in nitrogen and must be cleared from the bloodstream, such as urea, uric acid, and creatinine. A urinalysis can detect nitrogenous wastes of the mammalian body.
Urine plays an important role in the earth's nitrogen cycle. In balanced ecosystems, urine fertilizes the soil and thus helps plants to grow. Therefore, urine can be used as a fertilizer. Some animals use it to mark their territories. Historically, aged or fermented urine (known as lant) was also used for gunpowder production, household cleaning, tanning of leather and dyeing of textiles.
Human urine and feces are collectively referred to as human waste or human excreta, and are managed via sanitation systems. Livestock urine and feces also require proper management if the livestock population density is high.
Physiology
Most animals have excretory systems for elimination of soluble toxic wastes. In humans, soluble wastes are excreted primarily by the urinary system and, to a lesser extent in terms of urea, removed by perspiration. In placental mammals, the urinary system consists of the kidneys, ureters, urinary bladder, and urethra. The system produces urine by a process of filtration, reabsorption, and tubular secretion. The kidneys extract the soluble wastes from the bloodstream, as well as excess water, sugars, and a variety of other compounds. The resulting urine contains high concentrations of urea and other substances, including toxins. Urine flows from the kidneys through the ureter, bladder, and finally the urethra before passing through the urinary meatus.
Duration
Research looking at the duration of urination in a range of mammal species found that nine larger species urinated for 21 ± 13 seconds irrespective of body size. Smaller species, including rodents and bats, cannot produce steady streams of urine and instead urinate with a series of drops.
Characteristics
Quantity
Average urine production in adult humans is around of urine per person per day with a normal range of per person per day, produced in around 6 to 8 urinations per day depending on state of hydration, activity level, environmental factors, weight, and the individual's health. Producing too much or too little urine needs medical attention. Polyuria is a condition of excessive production of urine (> 2.5 L/day), oliguria when < 400 mL are produced, and anuria being < 100 mL per day.
Constituents
About 91–96% of urine consists of water. The remainder can be broadly characterized into inorganic salts, urea, organic compounds, and organic ammonium salts. Urine also contains proteins, hormones, and a wide range of metabolites, varying by what is introduced into the body.
The total solids in urine are on average per day per person. Urea is the largest constituent of the solids, constituting more than 50% of the total. The daily volume and composition of urine varies per person based on the amount of physical exertion, environmental conditions, as well as water, salt, and protein intakes. In healthy persons, urine contains very little protein and an excess is suggestive of illness, as with sugar. Organic matter, in healthy persons, also is reported to at most 1.7 times more matter than minerals. However, any more than that is suggestive of illness.
However, it is important to note that lesser amounts and concentrations of other compounds and ions are often present in urination of humans.
Color
Urine varies in appearance, depending principally upon a body's level of hydration, interactions with drugs, compounds and pigments or dyes found in food, or diseases. Normally, urine is a transparent solution ranging from colorless to amber, but is usually a pale yellow. Usually urination color comes primarily from the presence of urobilin. Urobilin is a final waste product resulting from the breakdown of heme from hemoglobin during the destruction of aging blood cells.
Colorless urine indicates over-hydration. Colorless urine in drug tests can suggest an attempt to avoid detection of illicit drugs in the bloodstream through over-hydration.
Bloody urine is termed hematuria, a symptom of a wide variety of medical conditions.
Reddish or brown urine may be caused by porphyria (not to be confused with the harmless, temporary pink or reddish tint caused by beeturia).
Pinkish urine can result from the consumption of beets (beeturia)
Dark yellow urine is often indicative of dehydration.
Orange urine due to certain medications such as rifampin and phenazopyridine
Dark orange to brown urine can be a symptom of jaundice, rhabdomyolysis, or Gilbert's syndrome.
Greenish urine can result from the consumption of asparagus or foods, beverages with green pigments, or from a urinary tract infection.
Blue urine can be caused by the ingestion of methylene blue (e.g., in medications) or foods or beverages with blue dyes.
Blue urine stains can be caused by blue diaper syndrome.
Purple urine may be due to purple urine bag syndrome.
Black or dark-colored urine is referred to as melanuria and may be caused by a melanoma or non-melanin acute intermittent porphyria.
Odor
Sometime after leaving the body, urine may acquire a strong "fish-like" odor because of contamination with bacteria that break down urea into ammonia. This odor is not present in fresh urine of healthy individuals; its presence may be a sign of a urinary tract infection.
The odor of normal human urine can reflect what has been consumed or specific diseases. For example, an individual with diabetes mellitus may present a sweetened urine odor. This can be due to kidney diseases as well, such as kidney stones. Additionally, the presence of amino acids in urine (diagnosed as maple syrup urine disease) can cause it to smell of maple syrup.
Eating asparagus can cause a strong odor reminiscent of the vegetable caused by the body's breakdown of asparagusic acid. Likewise consumption of saffron, alcohol, coffee, tuna fish, and onion can result in telltale scents. Particularly spicy foods can have a similar effect, as their compounds pass through the kidneys without being fully broken down before exiting the body.
pH
The pH normally is within the range of 5.5 to 7 with an average of 6.2. In persons with hyperuricosuria, acidic urine can contribute to the formation of stones of uric acid in the kidneys, ureters, or bladder. Urine pH can be monitored by a physician or at home.
A diet which is high in protein from meat and dairy, as well as alcohol consumption can reduce urine pH, whilst potassium and organic acids, such as from diets high in fruit and vegetables, can increase the pH and make it more alkaline.
Cranberries, popularly thought to decrease the pH of urine, have actually been shown not to acidify urine. Drugs that can decrease urine pH include ammonium chloride, chlorothiazide diuretics, and methenamine mandelate.
Density
Human urine has a specific gravity of 1.003–1.035.
Bacteria and pathogens
Urine is not sterile, not even in the bladder. In the urethra, epithelial cells lining the urethra are colonized by facultatively anaerobic Gram-negative rod and cocci bacteria. One study conducted in Nigeria isolated a total of 77 distinct bacterial strains from 100 healthy children (ages 5–11) as well as 39 strains from 33 cow urine samples, a considerable amount being pathogens. Pathogens identified and their percentages were:
The study also states:
Examination for medical purposes
Many physicians in ancient history resorted to the inspection and examination of the urine of their patients. Hermogenes wrote about the color and other attributes of urine as indicators of certain diseases. Abdul Malik Ibn Habib of Andalusia ( 862 AD) mentions numerous reports of urine examination throughout the Umayyad empire. Diabetes mellitus got its name because the urine is plentiful and sweet. The name uroscopy refers to any visual examination of the urine, including microscopy, although it often refers to the aforementioned prescientific or Proto-scientific forms of urine examination. Clinical urine tests today duly note the color, turbidity, and odor of urine but also include urinalysis, which chemically analyzes the urine and quantifies its constituents. A culture of the urine is performed when a urinary tract infection is suspected, as bacteriuria without symptoms does not require treatment. A microscopic examination of the urine may be helpful to identify organic or inorganic substrates and help in the diagnosis.
The color and volume of urine can be reliable indicators of hydration level. Clear and copious urine is generally a sign of adequate hydration. Dark urine is a sign of dehydration. The exception occurs when diuretics are consumed, in which case urine can be clear and copious and the person still be dehydrated.
Uses
Source of medications
Urine contains proteins and other substances that are useful for medical therapy and are ingredients in many prescription drugs (e.g., Ureacin, Urecholine, Urowave). Urine from postmenopausal women is rich in gonadotropins that can yield follicle stimulating hormone and luteinizing hormone for fertility therapy. One such commercial product is Pergonal.
Urine from pregnant women contains enough human chorionic gonadotropins for commercial extraction and purification to produce hCG medication. Pregnant mare urine is the source of estrogens, namely Premarin. Urine also contains antibodies, which can be used in diagnostic antibody tests for a range of pathogens, including HIV-1.
Urine can also be used to produce urokinase, which is used clinically as a thrombolytic agent.
Fertilizer
Cleaning
Given that urea in urine breaks down into ammonia, urine has been used for cleaning. In pre-industrial times, urine was used – in the form of lant or aged urine – as a cleaning fluid. Urine was also used for whitening teeth in Ancient Rome.
Gunpowder
Urine was used before the development of a chemical industry in the manufacture of gunpowder. Urine, a nitrogen source, was used to moisten straw or other organic material, which was kept moist and allowed to rot for several months to over a year. The resulting salts were washed from the heap with water, which was evaporated to allow collection of crude saltpeter crystals, that were usually refined before being used in making gunpowder.
Survival uses
The US Army Field Manual advises drinking urine for survival. The manual explains that drinking urine tends to worsen rather than relieve dehydration due to the salts in it, and that urine should not be consumed in a survival situation, even when there is no other fluid available. In hot weather survival situations, where other sources of water are not available, soaking cloth (a shirt for example) in urine and putting it on the head can help cool the body.
During World War I, Germans experimented with numerous poisonous gases as weapons. After the first German chlorine gas attacks, Allied troops were supplied with masks of cotton pads that had been soaked in urine. It was believed that the ammonia in the pad neutralized the chlorine. These pads were held over the face until the soldiers could escape from the poisonous fumes.
Urban legend states that urine works well against jellyfish stings. This scenario has appeared many times in popular culture including in the Friends episode "The One With the Jellyfish", an early episode of Survivor, as well as the films The Real Cancun (2003), The Heartbreak Kid (2007) and The Paperboy (2012). However, at best it is ineffective, and in some cases this treatment may make the injury worse.
Textiles
Urine has often been used as a mordant to help prepare textiles, especially wool, for dyeing. In the Scottish Highlands and Hebrides, the process of "waulking" (fulling) woven wool is preceded by soaking in urine, preferably infantile.
Olfactory communication
Urine plays a role in olfactory communication, since it contains semiochemicals that act as pheromones. The urine of predator species often contains kairomones that serve as a repellent against their prey species.
History
The fermentation of urine by bacteria produces a solution of ammonia; hence fermented urine was used in Classical Antiquity to wash cloth and clothing, to remove hair from hides in preparation for tanning, to serve as a mordant in dying cloth, and to remove rust from iron. Ancient Romans used fermented human urine (in the form of lant) to cleanse grease stains from clothing. The emperor Nero instituted a tax () on the urine industry, continued by his successor, Vespasian. The Latin saying ('money does not smell') is attributed to Vespasian – said to have been his reply to a complaint from his son about the unpleasant nature of the tax. Vespasian's name is still attached to public urinals in France (), Italy (), and Romania ().
Alchemists spent much time trying to extract gold from urine, which led to discoveries such as white phosphorus by German alchemist Hennig Brand when distilling fermented urine in 1669. In 1773 the French chemist Hilaire Rouelle discovered the organic compound urea by boiling urine dry.
Language
The English word urine (, ) comes from the Latin (-ae, f.), which is cognate with ancient words in various Indo-European languages that concern water, liquid, diving, rain, and urination (for example Sanskrit meaning 'it rains' or meaning 'water' and Greek meaning 'to urinate'). The onomatopoetic term piss predates the word urine, but is now considered vulgar. Urinate was at first used mostly in medical contexts. Piss is also used in such colloquialisms as to piss off, piss poor, and the slang expression pissing down to mean heavy rain. Euphemisms and expressions used between parents and children (such as wee, pee, number one and many others) have long existed.
Lant is a word for aged urine, originating from the Old English word referring to urine in general.
| Biology and health sciences | Urinary system | Biology |
2114529 | https://en.wikipedia.org/wiki/Atat%C3%BCrk%20Dam | Atatürk Dam | The Atatürk Dam (), originally the Karababa Dam, is the third largest dam in the world and it is a zoned rock-fill dam with a central core on the Euphrates River on the border of Adıyaman Province and Şanlıurfa Province in the Southeastern Anatolia Region of Turkey. Built both to generate electricity and to irrigate the plains in the region, it was renamed in honour of Mustafa Kemal Atatürk (1881–1938), the founder of the Turkish Republic. The construction began in 1983 and was completed in 1990. The dam and the hydroelectric power plant, which went into service after the upfilling of the reservoir was completed in 1992, are operated by the State Hydraulic Works (DSİ). The reservoir created behind the dam, called Atatürk Reservoir (), is the third largest in Turkey.
The dam is situated northwest of Bozova, Şanlıurfa Province, on state road D-875 from Bozova to Adıyaman. Centerpiece of the 22 dams on the Euphrates and the Tigris, which comprise the integrated, multi-sector, Southeastern Anatolia Project (, known as GAP), it is one of the world's largest dams. The Atatürk Dam, one of the five operational dams on the Euphrates as of 2008, was preceded by Keban and Karakaya dams upstream and followed by Birecik and the Karkamış dams downstream. Two more dams on the river have been under construction.
The dam embankment is and . The hydroelectric power plant (HEPP) has a total installed power capacity of 2,400 MW and generates 8,900 GW·h electricity annually.
The total cost of the dam project was about .
The dam was depicted on the reverse of the Turkish one-million-lira banknotes of 1995–2005 and of the 1 new lira banknote of 2005–2009.
Dam
The initial development project for the southeastern region of Turkey was presented in 1970. As the objectives for regional development have changed significantly and the ambitions have grown in the 1970s, the original plan underwent major modifications. The most important change in the project was abandoning the Middle Karababa Dam design, and adopting the design of the Atatürk Dam to increase the storage and power generation capacities of the dam.
Dolsar Engineering and ATA Construction, two prominent Turkish companies, signed for the building of the dam. The construction of the cofferdam began in 1985 and was completed in 1987. The fill work for the main dam lasted from 1987 to 1990. The Atatürk Dam, listed in international construction publications as the world's largest construction site, was completed in a world record time of around 50 months.
The rock-fill dam undergoes deformations that are regularly and systematically monitored since 1990 with different types of sensors. It is estimated that the central portion of the dam crest has settled by around since the end of the construction. Settlement of the dam crest up to has been measured since the start of the detailed geodetic monitoring in 1992. The maximum horizontal (radial) deformation measured is about .
The permeation grouting work was carried out by subcontractor Solétanche Bachy and the rehabilitation work for the post-tensioning of the dam crest with ground anchors by Vorspann System Losinger International (VSL).
Hydroelectric power plant
The HEPP of the Atatürk Dam is the biggest of a series of 19 power plants of the GAP project. It consists of eight Francis turbine and generator groups of 300 MW each, supplied by Sulzer Escher Wyss and ABB Asea Brown Boveri respectively. The up to steel pressure pipes (penstocks) with a total weight of 26.600 tons were supplied and installed by the German NOELL company (today DSD NOELL). The power plant's first two power units came on line in 1992 and it became fully operational in December 1993. The HEPP can generate 8,900 GWh of electricity annually. Its capacity makes up around one third of the total capacity of the GAP project.
During the periods of low demand for electricity, only one of the eight units of the HEPP is in operation while in times of high demand, all the eight units are in operation. Hence, depending upon the energy demand and the state of the interconnected system, the amount of water to be released from the HEPP might vary between 200 and 2,000 m3/s in one day.
Irrigation
Originating in the mountains of eastern Anatolia and flowing southwards to Syria and Iraq, the Euphrates and the Tigris are very irregular rivers, used to cause great problems each year with droughts in summer and flooding in winter. The water of the Euphrates River is regulated by means of large reservoirs of the Keban and Atatürk Dams. However, the waters released from the HEPPs of those dams also need to be regulated. The Birecik and the Karkamış Dams downstream the Atatürk Dam are constructed for the purpose of harnessing the waters released from large-scale dams and HEPPs.
Nearly of arable land in the Şanlıurfa-Harran and Mardin-Ceylanpınar plains in upper Mesopotamia is being irrigated via gravity-flow with water diverted from the Atatürk Dam through the Şanlıurfa Tunnels system, which consists of two parallel tunnels, each long and in diameter. The flow rate of water through the tunnels is about , which makes one-third of the total flow of the Euphrates. The tunnels are the largest in the world, in terms of length and flow rate, built for irrigation purposes. The first tunnel was completed in 1995 and the other in 1996. The reservoir behind the dam will irrigate another 406,000 ha by pumping for a total of 882,000 ha.
The Atatürk Dam and the Şanlıurfa Tunnel system are two major components of the GAP project. Irrigation started in the Harran Plain in the spring of 1995. The impact of the irrigation on the economy of the region is significant. In ninety percent of the irrigated area, cotton is planted. Irrigation expansion within the Harran plains also increased Southeastern Anatolia's cotton production from 164,000 to 400,000 metric tons in 2001, or nearly sixty percent. With almost 50% share of the country's cotton production, the region developed to the leader in Turkey.
Reservoir lake
The Atatürk Reservoir, extending over an area of with a water volume of 48.7 km3 (63,400 million cu yd), ranks third in size in Turkey after Lake Van and Lake Tuz. The reservoir water level touched amsl in 1994. Since then, it varies between 526 and 537 m amsl. The full reservoir level is , and the minimum operation level is amsl.
Some 10 towns and 156 villages of three provinces are located around the Lake Atatürk Dam. The lake provides a fisheries and recreation site. For transportation purposes, several ferries have been operated in the reservoir. The reservoir lake is called "sea" by local people.
Geostrategic importance
About 90% of Euphrates' total annual flow originates in Turkey, while the remaining part is added in Syria, but nothing is contributed further downstream in Iraq. In general, the stream varies greatly in its flow from season to season and year to year. As an example, the annual flow at the border with Syria ranged from in 1961 to in 1963.
One of the most important legal texts on the waters of the Euphrates-Tigris river system is the protocol annexed to the 1946 Treaty of Friendship and Good Neighborly Relations between Iraq and Turkey. The protocol provided the control and management of the Euphrates and the Tigris depending to a large extent on the regulations of flow in Turkish source areas. Turkey agreed to begin monitoring the two border-crossing rivers and to share related data with Iraq. In 1980, Turkey and Iraq further specified the nature of the earlier protocol by forming a joint committee on technical issues, which Syria joined later in 1982 as well. Turkey unilaterally guaranteed to allow 15.75 km3/year (500 m3/s) of water across the border to Syria without any formal agreement on the sharing of the Euphrates water.
Mid-January 1990, when the first phase of the dam was completed, Turkey held back the flow of the Euphrates entirely for a month to begin filling up the reservoir. Turkey had notified Syria and Iraq by November 1989 of her decision to fill the reservoir over a period of one month explaining the technical reasons and providing a detailed program for making up for the losses. The downstream neighbors protested vehemently. At this point, the Atatürk Dam has cut the flow from the Euphrates by about a third.
Syria and Iraq claim to be suffering severe water shortages due to the GAP development. Both countries allege that Turkey is intentionally withholding supplies from its downstream neighbors, turning water into a weapon. Turkey denies these claims, and insists it has always supplied its southern neighbors with the promised minimum of . It argues that Iraq and Syria in fact benefit from the regulated water by the dams as they protect all three riparian countries from seasonal droughts and floods.
Two damaging earthquakes of M w 5.5 and M w 5.1 occurred in the town of Samsat near the Atatürk Reservoir in 2017 and 2018, respectively. The spatio-temporal evolution of seismicity and its source properties in relation to the temporal water-level variations and the stresses resulting from surface loading and pore-pressure diffusion shows the water-level and seismicity rate are anti-correlated in this dam, which is explained by the stabilization effect of the gravitational induced stress imposed by water loading on the local faults. The overall effective stress in the seismogenic zone increased over decades due to pore-pressure diffusion, explaining the enhanced background seismicity during recent years.
| Technology | Dams | null |
2116016 | https://en.wikipedia.org/wiki/CIE%201931%20color%20space | CIE 1931 color space | In 1931 the International Commission on Illumination (CIE) published the CIE 1931 color spaces which define the relationship between the visible spectrum and the visual sensation of specific colors by human color vision. The CIE color spaces are mathematical models that create a "standard observer", which attempts to predict the perception of unique hues of color. These color spaces are essential tools that provide the foundation for measuring color for industry, including inks, dyes, and paints, illumination, color imaging, etc. The CIE color spaces contributed to the development of color television, the creation of instruments for maintaining consistent color in manufacturing processes, and other methods of color management.
The initials CIE come from the French name "Commission Internationale de l'éclairage" (International Commission on Illumination), which has maintained and developed many of the standards in use today relating to colorimetry. The CIE color spaces were created using data from a series of experiments, where human test subjects adjusted red, green, and blue primary colors to find a visual match to a second, pure color. The original experiments were conducted in the mid 1920s by using ten observers and John Guild using seven observers. The experimental results were combined, creating the CIE RGB color space. The CIE XYZ color space was derived from CIE RGB in an effort to simplify the math.
The CIE 1931 XYZ color space is still widely used, even though it is not perceptually uniform in relation to human vision. In 1976 the CIE published the CIELUV and CIELAB color spaces, which are derived from XYZ, and are intended to provide more uniform predictions relative to human perception.
Background: the human eye and tristimulus values
The human eye with normal vision has three kinds of cone cells that sense light, having peaks of spectral sensitivity in short ("S", – ), medium ("M", – ), and long ("L", – ) wavelengths. These cone cells underlie human color perception in conditions of medium and high brightness; in very dim light color vision diminishes, and the low-brightness, monochromatic "night vision" receptors, denominated "rod cells", become effective. Thus, three parameters corresponding to levels of stimulus of the three kinds of cone cells, in principle describe any human color sensation. Weighting a total light power spectrum by the individual spectral sensitivities of the three kinds of cone cells renders three effective values of stimulus; these three values compose a tristimulus specification of the objective color of the light spectrum. The three parameters, denoted "S", "M", and "L", are indicated using a 3-dimensional space denominated the "LMS color space", which is one of many color spaces devised to quantify human color vision.
A color space maps a range of physically produced colors from mixed light, pigments, etc. to an objective description of color sensations registered in the human eye, typically in terms of tristimulus values, but not usually in the LMS color space defined by the spectral sensitivities of the cone cells. The tristimulus values associated with a color space can be conceptualized as amounts of three primary colors in a tri-chromatic, additive color model. In some color spaces, including the LMS and XYZ spaces, the primary colors used are not real colors in the sense that they cannot be generated in any light spectrum.
The CIE XYZ color space encompasses all color sensations that are visible to a person with average eyesight. That is why CIE XYZ tristimulus values are a device-invariant representation of color. It serves as a standard reference against which many other color spaces are defined. A set of color-matching functions, like the spectral sensitivity curves of the LMS color space, but not restricted to non-negative sensitivities, associates physically produced light spectra with specific tristimulus values.
Consider two light sources composed of different mixtures of various wavelengths. Such light sources may appear to be the same color; this effect is called "metamerism." Such light sources have the same apparent color to an observer when they produce the same tristimulus values, regardless of the spectral power distributions of the sources.
Most wavelengths stimulate two or all three kinds of cone cell because the spectral sensitivity curves of the three kinds overlap. Certain tristimulus values are thus physically impossible: e.g. LMS tristimulus values that are non-zero for the M component and zero for both the L and S components. Furthermore pure spectral colors would, in any normal trichromatic additive color space, e.g., the RGB color spaces, imply negative values for at least one of the three primaries because the chromaticity would be outside the color triangle defined by the primary colors. To avoid these negative RGB values, and to have one component that describes the perceived brightness, "imaginary" primary colors and corresponding color-matching functions were formulated. The CIE 1931 color space defines the resulting tristimulus values, in which they are denoted by "X", "Y", and "Z". In XYZ space, all combinations of non-negative coordinates are meaningful, but many, such as the primary locations [1, 0, 0], [0, 1, 0], and [0, 0, 1], correspond to imaginary colors outside the space of possible LMS coordinates; imaginary colors do not correspond to any spectral distribution of wavelengths and therefore have no physical reality.
Meaning of X, Y and Z
In the CIE 1931 model, Y is the luminance, Z is quasi-equal to blue (of CIE RGB), and X is a mix of the three CIE RGB curves chosen to be nonnegative (see ). Setting Y as luminance has the useful result that for any given Y value, the XZ plane will contain all possible chromaticities at that luminance.
The unit of the tristimulus values , , and is often arbitrarily chosen so that or is the brightest white that a color display supports. In this case, the Y value is known as the relative luminance. The corresponding whitepoint values for and can then be inferred using the standard illuminants.
Since the XYZ values are defined much earlier than the characterization of cone cells in the 1950s (by Ragnar Granit), the physiological meaning of these values are known only much later. The Hunt-Pointer-Estevez matrix from the 1980s relates XYZ with LMS. When inverted, it shows how the three cone responses add up to XYZ functions:
In other words, the Z value is solely made up of the S cone response, the Y value a mix of L and M responses, and X value a mix of all three. This fact makes XYZ values analogous to, but different from, the LMS cone responses of the human eye.
CIE standard observer
Due to the distribution of cones in the eye, the tristimulus values depend on the observer's field of view. To eliminate this variable, the CIE defined a color-mapping function called the standard (colorimetric) observer, to represent an average human's chromatic response within a 2° arc inside the fovea. This angle was chosen owing to the belief that the color-sensitive cones resided within a 2° arc of the fovea. Thus the CIE 1931 Standard Observer function is also known as the CIE 1931 2° Standard Observer. A more modern but less-used alternative is the CIE 1964 10° Standard Observer, which is derived from the work of Stiles and Burch, and Speranskaya.
For the 10° experiments, the observers were instructed to ignore the central 2° spot. The 1964 Supplementary Standard Observer function is recommended when dealing with more than about a 4° field of view. Both standard observer functions are discretized at wavelength intervals from to and distributed by the CIE. All corresponding values have been calculated from experimentally obtained data using interpolation. The standard observer is characterized by three color matching functions.
There is also a -interval dataset of CIE 1931 and CIE 1964 provided by Wyszecki 1982. A CIE publication in 1986 appears also to have a 1 nm dataset, probably using the same data. Like the regular dataset, this dataset is also derived from interpolation.
The derivation of the CIE standard observer from color matching experiments is given below, after the description of the CIE RGB space.
Color matching functions
The CIE's color matching functions , and are the numerical description of the chromatic response of the observer (described above). They can be thought of as the spectral sensitivity curves of three linear light detectors yielding the CIE tristimulus values X, Y and Z. Collectively, these three functions describe the CIE standard observer.
Analytical approximation
Table lookup can become impractical for some computational tasks. Instead of referring to the published table, the CIE XYZ color matching functions can be approximated by a sum of Gaussian functions, as follows:
Let g(x) denote a piecewise-Gaussian function, defined by
That is, g(x) resembles a bell curve with its peak at , a spread/standard deviation of to the left of the mean, and spread of to the right of the mean. With the wavelength λ measured in nanometers, we then approximate the 1931 color matching functions:
The squared differences between the above approximation and the measured CIE xyz color matching functions is less than the within-observer variance encountered in the experimental measurements used to form the CIE standards. It is also possible to use fewer gaussian functions, with one gaussian for each "lobe". CIE 1964 fits well with a one-lobe function.
The CIE XYZ color matching functions are nonnegative, and lead to nonnegative XYZ coordinates for all real colors (that is, for nonnegative light spectra). Other observers, such as for the CIE RGB space or other RGB color spaces, are defined by other sets of three color-matching functions, not generally nonnegative, and lead to tristimulus values in those other spaces, which may include negative coordinates for some real colors.
Computing XYZ from spectral data
Emissive case
The tristimulus values for a color with a spectral radiance Le,Ω,λ are given in terms of the standard observer by:
where is the wavelength of the equivalent monochromatic light (measured in nanometers), and customary limits of the integral are .
The values of X, Y, and Z are bounded if the radiance spectrum Le,Ω,λ is bounded.
Reflective and transmissive cases
The reflective and transmissive cases are very similar to the emissive case, with a few differences. The spectral radiance Le,Ω,λ is replaced by the spectral reflectance (or transmittance) S(λ) of the object being measured, multiplied by the spectral power distribution of the illuminant I(λ).
where
K is a scaling factor (usually 1 or 100), and is the wavelength of the equivalent monochromatic light (measured in nanometers), and the standard limits of the integral are .
CIE xy and the CIE xyY color space
Since the human eye has three types of color sensors that respond to different ranges of wavelengths, a full plot of all visible colors is a three-dimensional figure. However, the concept of color can be divided into two parts: brightness and chromaticity. For example, the color white is a bright color, while the color grey is considered to be a less bright version of that same white. In other words, the chromaticity of white and grey are the same while their brightness differs.
The CIE XYZ color space was deliberately designed so that the Y parameter is also a measure of the luminance of a color. The chromaticity is then specified by the two derived parameters x and y, two of the three normalized values being functions of all three tristimulus values X, Y, and Z:
That is, because each tristimulus parameter, X, Y, Z, is divided by the sum of all three, the resulting values, x, y, z, each represent a proportion of the whole and so their sum must be equal to one. Therefore, the value z can be deduced by knowing x and y, and consequently the latter two values are sufficient for describing the chromaticity of any color.
The derived color space specified by x, y, and Y is known as the CIE xyY color space and is widely used to specify colors in practice.
The X and Z tristimulus values can be calculated back from the chromaticity values x and y and the Y tristimulus value:
The figure on the right shows the related chromaticity diagram. The outer curved boundary is the spectral locus, with wavelengths shown in nanometers. The chromaticity diagram is a tool to specify how the human eye will experience light with a given spectrum. It cannot specify colors of objects (or printing inks), since the chromaticity observed while looking at an object depends on the light source as well.
Mathematically the colors of the chromaticity diagram occupy a region of the real projective plane.
The chromaticity diagram illustrates a number of interesting properties of the CIE XYZ color space:
The diagram represents all of the chromaticities visible to the average person. These are shown in color and this region is called the gamut of human vision. The gamut of all visible chromaticities on the CIE plot is the tongue-shaped or horseshoe-shaped figure shown in color. The curved edge of the gamut is called the spectral locus and corresponds to monochromatic light (each point representing a pure hue of a single wavelength), with wavelengths listed in nanometers. The straight edge on the lower part of the gamut is called the line of purples. These colors, although they are on the border of the gamut, have no counterpart in monochromatic light. Less saturated colors appear in the interior of the figure with white at the center.
It is seen that all visible chromaticities correspond to non-negative values of x, y, and z (and therefore to non-negative values of X, Y, and Z).
If one chooses any two points of color on the chromaticity diagram, then all the colors that lie in a straight line between the two points can be formed by mixing these two colors. It follows that the gamut of colors must be convex in shape. All colors that can be formed by mixing three sources are found inside the triangle formed by the source points on the chromaticity diagram (and so on for multiple sources).
An equal mixture of two equally bright colors will not generally lie on the midpoint of that line segment. In more general terms, a distance on the CIE xy chromaticity diagram does not correspond to the degree of difference between two colors. In the early 1940s, David MacAdam studied the nature of visual sensitivity to color differences, and summarized his results in the concept of a MacAdam ellipse. Based on the work of MacAdam, the CIE 1960, CIE 1964, and CIE 1976 color spaces were developed, with the goal of achieving perceptual uniformity (have an equal distance in the color space correspond to equal differences in color). Although they were a distinct improvement over the CIE 1931 system, they were not completely free of distortion.
It can be seen that, given three real sources, these sources cannot cover the gamut of human vision. Geometrically stated, there are no three points within the gamut that form a triangle that includes the entire gamut; or more simply, the gamut of human vision is not a triangle.
Light with a flat power spectrum in terms of wavelength (equal power in every interval) corresponds to the point .
Mixing colors specified with the CIE xy chromaticity diagram
When two or more colors are additively mixed, the x and y chromaticity coordinates of the resulting color (xmix,ymix) may be calculated from the chromaticities of the mixture components (x1,y1; x2,y2; …; xn,yn) and their corresponding luminances (L1, L2, …, Ln) with the following formulas:
These formulas can be derived from the previously presented definitions of x and y chromaticity coordinates by taking advantage of the fact that the tristimulus values X, Y, and Z of the individual mixture components are directly additive. In place of the luminance values (L1, L2, etc.) one can alternatively use any other photometric quantity that is directly proportional to the tristimulus value Y (naturally meaning that Y itself can also be used as well).
As already mentioned, when two colors are mixed, the resulting color xmix, ymix will lie on the straight line segment that connects these colors on the CIE xy chromaticity diagram. To calculate the mixing ratio of the component colors x1,y1 and x2,y2 that results in a certain xmix,ymix on this line segment, one can use the formula
where L1 is the luminance of color x1,y1 and L2 the luminance of color x2,y2. Because ymix is unambiguously determined by xmix and vice versa, knowing just one or the other of them is enough for calculating the mixing ratio. In accordance with the remarks concerning the formulas for xmix and ymix, the mixing ratio L1/L2 may well be expressed in terms of other photometric quantities than luminance.
Definition of the CIE XYZ color space
CIE RGB color space
The first step in developing the CIE XYZ color space is the measurement of the CIE RGB color space. The CIE RGB color space is one of many RGB color spaces, distinguished by a particular set of monochromatic (single-wavelength) primary colors.
In the 1920s, two independent experiments on human color perception were conducted by W. David Wright with ten observers, and John Guild with seven observers. Their results laid the foundation for the trichromatic CIE XYZ color space specification.
The experiments were conducted by using a circular split screen (a bipartite field) 2 degrees in diameter, which is the angular size of the human fovea. On one side a test color was projected while on the other an observer-adjustable color was projected. The adjustable color was a mixture of the three monochromatic primary colors, each with adjustable brightness. The observer would alter the brightness of each of the three primary beams until a match to the test color was observed.
If the test color were simply a monochromatic color at wavelength λ, and if it could be matched by a combination of the three primaries at relative intensities , , and respectively, then a tabulation of these values at various λ will estimate three functions of wavelength. These are the RGB color-matching functions. Any spectral distribution can be thought of as a combination of a number of monochromatic sources at varying intensities, so that (by Grassmann's laws) integrating the color matching functions with that spectral distribution will yield the intensities of the three primaries necessary to match it. The problem is that the three primaries can only produce colors which lie within their gamut - the triangle in color space formed by the primaries, which never touches the monochromatic locus nor the purple line except at the three primaries. In other words, there is no monochromatic source that can be matched by a combination of the three primaries, except at the wavelengths of the three primaries themselves. However, by adding one of the primaries to the monochromatic test color, the test color can be brought into the RGB gamut, allowing a match to be made. Adding a primary to the monochromatic test color is effectively the same as subtracting it from the adjustable color, which of course cannot be done since it is impossible to have a negative intensity for any of the primaries.
For wavelengths between the blue and green primaries, some red primary must be added to allow matching, resulting in negative values of . Likewise, between the green and red primaries, some blue must be added and will be negative. For wavelengths below the wavelength of the blue primary, or above the wavelength of the red primary, some green must be added and will be negative. In each case, the remaining two color matching functions will be positive. It can be seen that the deviation of the RGB gamut from the complete gamut is rather small except between the blue and green primaries at 435.8 and 546.1 nm. In this wavelength band, rather large amounts of the red primary needed to be added to the test color, and it is in this band that the red color matching function has rather large negative values. In their regions of negative values, the green and blue matching functions have rather small negative values.
Although Wright and Guild's experiments were carried out using various primaries at various intensities, and although they used a number of different observers, all of their results were summarized by the standardized CIE RGB color matching functions , , and , obtained using three monochromatic primaries at standardized wavelengths of (red), (green) and (blue). The (un-normalized) color matching functions are the amounts of primaries needed to match the monochromatic test primary. These functions are shown in the plot on the right (CIE 1931). and are zero at , and are zero at and and are zero at , since in these cases the test color is one of the primaries. The primaries with wavelengths and were chosen because they are easily reproducible monochromatic lines of a mercury vapor discharge. The wavelength, which in 1931 was difficult to reproduce as a monochromatic beam, was chosen because the eye's perception of color is rather unchanging at this wavelength, and therefore small errors in wavelength of this primary would have little effect on the results.
The color matching functions and primaries were settled upon by a CIE special commission after considerable deliberation. The cut-offs at the short- and long-wavelength side of the diagram are chosen somewhat arbitrarily; the human eye can actually see light with wavelengths up to about , but with a sensitivity that is many thousand times lower than for green light. These color matching functions define what is known as the "1931 CIE standard observer". Rather than specify the brightness of each primary, the curves are normalized to have constant area beneath them. This area is fixed to a particular value by specifying that
The resulting normalized color matching functions are then scaled in the r:g:b ratio of 1:4.5907:0.0601 for source luminance and 72.0962:1.3791:1 for source radiance to reproduce the true color matching functions. By proposing that the primaries be standardized, the CIE established an international system of objective color notation.
Given these scaled color matching functions, the RGB tristimulus values for
a color with a spectral power distribution would then be given by:
These are all inner products and can be thought of as a projection of an infinite-dimensional spectrum to a three-dimensional color.
Grassmann's Laws
One might ask: "Why is it possible that Wright and Guild's results can be summarized using different primaries and different intensities from those actually used?" One might also ask: "What about the case when the test colors being matched are not monochromatic?" The answer to both of these questions lies in the (near) linearity of human color perception. This linearity is expressed in Grassmann's laws of color.
The CIE RGB space can be used to define chromaticity in the usual way: The chromaticity coordinates are r, g and b where:
Construction of the CIE XYZ color space from the Wright–Guild data
Having developed an RGB model of human vision using the CIE RGB matching functions, the members of the special commission wished to develop another color space that would relate to the CIE RGB color space. It was assumed that Grassmann's law held, and the new space would be related to the CIE RGB space by a linear transformation. The new space would be defined in terms of three new color matching functions
, , and as described above. The new color space would be chosen to have the following desirable properties:
The new color matching functions were to be everywhere greater than or equal to zero. In 1931, computations were done by hand or slide rule, and the specification of positive values was a useful computational simplification.
The color matching function would be exactly equal to the photopic luminous efficiency function V(λ) for the "CIE standard photopic observer". The luminance function describes the variation of perceived brightness with wavelength. The fact that the luminance function could be constructed by a linear combination of the RGB color matching functions was not guaranteed by any means but might be expected to be nearly true due to the near-linear nature of human sight. Again, the main reason for this requirement was computational simplification.
For the constant energy white point, it was required that .
By virtue of the definition of chromaticity and the requirement of positive values of x and y, it can be seen that the gamut of all colors will lie inside the triangle [1, 0], [0, 0], [0, 1]. It was required that the gamut fill this space practically completely.
It was found that the color matching function could be set to zero above while remaining within the bounds of experimental error. For computational simplicity, it was specified that this would be so.
In geometrical terms, choosing the new color space amounts to choosing a new triangle in rg chromaticity space. In the figure above-right, the rg chromaticity coordinates are shown on the two axes in black, along with the gamut of the 1931 standard observer. Shown in red are the CIE xy chromaticity axes which were determined by the above requirements. The requirement that the XYZ coordinates be non-negative means that the triangle formed by Cr, Cg, Cb must encompass the entire gamut of the standard observer. The line connecting Cr and Cb is fixed by the requirement that the function be equal to the luminance function. This line is the line of zero luminance, and is called the alychne. The requirement that the function be zero above means that the line connecting Cg and Cr must be tangent to the gamut in the region of Kr. This defines the location of point Cr. The requirement that the equal energy point be defined by puts a restriction on the line joining Cb and Cg, and finally, the requirement that the gamut fill the space puts a second restriction on this line to be very close to the gamut in the green region, which specifies the location of Cg and Cb. The above described transformation is a linear transformation from the CIE RGB space to XYZ space. The standardized transformation settled upon by the CIE special commission was as follows:
The numbers in the conversion matrix below are exact, with the number of digits specified in CIE standards.
The above matrix is balanced for the equi-energy stimulus: it has coordinates (1,1,1) in both RGB and XYZ coordinates.
While the above matrix is exactly specified in standards, the inverse is left unspecified so that it can be approximated to machine precision to reduce round-trip rounding errors. Its values can be computed precisely using rational numbers:
Which has these approximate values:
The XYZ primaries will have XYZ coordinates [1,0,0], [0,1,0], and [0,0,1] in XYZ space, so the columns of the inverse matrix above specify the XYZ primaries ( Cr, Cg and Cb) in RGB space. Dividing each column by its sum will give the coordinates of the XYZ primaries in rgb space which yields:
Cr = {1.27496, -0.27777, 0.00280576}
Cg = {-1.7393, 2.76726, -0.0279521}
Cb = {-0.743104, 0.140911, 1.60219}
The r and g coordinates of the XYZ primaries are indicated in the rg chromaticity space diagram above.
The integrals of the XYZ color matching functions must all be equal by requirement 3 above, and this is set by the integral of the photopic luminous efficiency function by requirement 2 above. The tabulated sensitivity curves have a certain amount of arbitrariness in them. The shapes of the individual X, Y and Z sensitivity curves can be measured with a reasonable accuracy. However, the overall luminosity curve (which in fact is a weighted sum of these three curves) is subjective, since it involves asking a test person whether two light sources have the same brightness, even if they are in completely different colors. Along the same lines, the relative magnitudes of the X, Y, and Z curves are arbitrary. Furthermore, one could define a valid color space with an X sensitivity curve that has twice the amplitude. This new color space would have a different shape. The sensitivity curves in the CIE 1931 and 1964 XYZ color spaces are scaled to have equal areas under the curves.
Subsequent refinements
A few other XYZ-style color-matching functions have been available, correcting for known issues in the original 1931 color space. These functions imply their own XYZ-like and xyY-like color spaces.
Judd and Vos corrections for the 2° CMF
The most serious problem with the 1931 CIE XYZ color matching functions is the error in the photopic Y (or function on the blue end of the spectrum. The Judd (1951) and its following Vos (1978) corrections sought to correct for the issue without deviating from the original methodology.
CIE 1964 X10Y10Z10
X10Y10Z10 (also written XYZ10 and analogously for the following) is the XYZ-style color space defined using the CIE 1964 10° observer CMFs. The 3 CMFs are mainly derived from Stiles and Burch's RGB color-matching functions, which unlike the Wright–Guild functions (and the subsequent Judd–Vos corrections) are "directly measured", freeing them from the reconstruction errors of the 1931 functions.
Stiles and Burch also published a set of 2° RGB color-matching functions; however, no XYZ space derived from them has been formally recognized by the CIE.
CIE 170-2 XFYFZF
XFYFZF is the XYZ-style color space defined using the Stockman & Sharpe (2000) physiological 2° observer, which is in turn a linear combination of the LMS cone response functions. The CMF data, along with the physiological 10° dataset, is available from the Colour & Vision Research laboratory of University College London down to 0.1 nm resolution.
CIE 170-2 XF,10YF,10ZF,10
This space is based on the Stockman & Sharpe (2000) physiological 10° observer.
According to Konica Minolta, the older CIE 1931 CMF exhibits metamerism failure (failure to predict when colors appear the same) for wide color gamut displays containing narrowband emitters like OLED, whereas the 2015 XYZF CMF is not affected. Older Sony manuals recommend using the Judd-Vos correction by applying an offset to the white point depending on the display technology used.
| Physical sciences | Basics | Physics |
2116156 | https://en.wikipedia.org/wiki/Pharaoh%20ant | Pharaoh ant | The pharaoh ant (Monomorium pharaonis) is a small (2 mm) yellow or light brown, almost transparent ant notorious for being a major indoor nuisance pest, especially in hospitals. A cryptogenic species, it has now been introduced to virtually every area of the world, including Europe, the Americas, Australasia and Southeast Asia. It is a major pest in the United States, Australia, and Europe. The ant's common name is possibly derived from the mistaken belief that it was one of the Egyptian (pharaonic) plagues.
This species is polygynouseach colony contains many queensleading to unique caste interactions and colony dynamics. This also allows the colony to fragment into bud colonies quickly.
Pharaoh ants are a tropical species, but they also thrive in buildings almost anywhere, even in temperate regions provided central heating is present.
Physical characteristics
Pharaoh workers are about 1.5– long. They are light yellow to reddish brown in color with a darker abdomen. Pharaoh ant workers have a non-functional stinger used to generate pheromones. The petiole (narrow waist between the thorax and abdomen) has two nodes and the thorax has no spines. Pharaoh ant eyesight is poor and they possess on average 32 ommatidia. The antennal segments end in a distinct club with three progressively longer segments.
Males are about long, black, winged (but do not fly). Queens are dark red and long. They initially have wings that are lost soon after mating, but do not fly.
Life cycle
The pharaoh ant queen can lay hundreds of eggs in her lifetime. Most lay 10 to 12 eggs per batch in the early days of egg production and only four to seven eggs per batch later. At 27 °C (80 °F) and 80 percent relative humidity, eggs hatch in five to seven days. The larval period is 18 to 19 days, pre-pupal period three days and pupal period nine days. About four more days are required to produce sexual female and male forms. From egg to sexual maturity, it takes the pharaoh ant about 38 to 45 days, depending on temperature and relative humidity. They breed continuously throughout the year in heated buildings and mating occurs in the nest. Mature colonies contain several queens, winged males, workers, eggs, larvae, pre-pupae and pupae.
Colony proliferation
Each colony produces sexually reproductive individuals roughly twice a year. However, colonies raised in a laboratory can be manipulated to produce sexuals at any time of year. Colonies proliferate by "budding" (also called "satelliting" or "fractionating"), where a subset of the colony including queens, workers and brood (eggs, larvae and pupae) leave the main colony for an alternative nest site.
Pharaoh ant colonies appear to prefer familiar nests to novel nests while budding. This suggests the ability for colonies to remember certain qualities of their living space. However, if the novel (unfamiliar) nest is of superior quality, the colony may initially move toward the familiar, but will eventually select the unfamiliar. The colony assumes the familiar nest is preferable, unless they sense better qualities in the novel nest. This decision-making process seeks to minimize the time the colony is without a nest while optimizing the nest the colony finally chooses.
The number of available budding locations has a large effect on colony fragmentation. A large number of bud nests results in small colony fragments, indicating that the colony has the ability to control size and caste ratios. However, a minimum group size of 469 individuals appears preferred by the species. Amount of fragmentation does not have an effect on food distribution. After budding, nest units do not compete for resources, but rather act cooperatively. This is evolutionarily explained by the high amount of genetic relatedness among these nest units. In addition, major disturbances to the central nest cause the colony to abandon it and flee to a bud nest. Thus, nest units may exchange individuals after budding occurs, further explaining their cooperative behavior.
In Australia, Monomorium species is particularly successful. This fact is particularly curious because of the presence of a very aggressive ant family, Iridomyrmex, which is quite proficient at interference competition. Iridomyrmex ants are able to quickly seek out food sources and prevent other ant species from reaching them. However, unlike other ant species, Monomorium species, despite their unaggressive nature and small size, are able to thrive even in areas where Iridomyrmex dominates. This success can be attributed to their efficient foraging strategy, and their novel use of venom alkaloids, repellant chemical signals. With these two behaviors, Monomorium species can rapidly monopolize and defend food sources.
Pheromones
Pharaoh ants utilize three types of pheromones. One is a long-lasting attractive chemical that is used to build a trail network. It remains detectable even if the ants do not use the trail for several days. Pharaoh ants cease activity at night and begin each day of work at around 8 am, yet parts of the trail network are identical each day. The second pheromone is also attractive, but will decay to imperceptible amounts in a matter of minutes without reapplication. This pheromone is useful in marking food sources, as these are unpredictable and the colony must be able to respond to environmental changes quickly. Individuals will not waste their time on an unprofitable trail route. The third pheromone is a repellant. Pharaoh ants were the first species found to use a negative trail pheromone. If an individual finds an unprofitable area with little food or significant danger, it will release this repellant pheromone, which will warn others and cause them to look elsewhere. While positive pheromones indicating lucrative foraging sites are very common in social insects, the pharaoh ant's negative pheromone is highly unusual. Like the food source marker, the negative pheromone is volatile, decaying roughly two hours after being emitted. It may even be insecticidal in some cases. It is so powerful that an individual can detect it from away. Pharaoh ants utilize this pheromone near forks in the trail network, and an ant that detects it will begin to walk in a zigzag manner.
Both the attractive and repellent pheromones are utilized in the decision-making an ant must make while foraging. The repellent pheromone is especially useful in the repositioning of trails after a new food source has been introduced. It also helps prevent ants from concentrating on an undesirable trail. Thus, the repellant pheromone makes the pharaoh ant a particularly efficient forager. Despite their extreme importance, there is an adaptive value to using pheromones sparingly, as it streamlines communication during important decision-making situations, such as a nest migration.
Foraging
Pharaoh ants use a positive feedback system of foraging. Each morning, scouts will search for food. When one finds it, it will immediately return to the nest. This causes several ants to follow the successful scout's trail back to the food source. Soon, a large group will be upon the food. Scouts are thought to use both chemical and visual cues to remain aware of the nest location and find their way. If the colony is exploring a new region, they employ a land rush tactic, in which a large number of foragers randomly search, constantly releasing pheromones.
Even though M. pharaonis is most often thought an indoor pest, foraging has been found to be more prevalent outside. Even inside colonies were found to forage close to windows, indicating a propensity for outdoor environment.
Trails
Even though scouts search independently, they use a distinct trail system that remains somewhat constant from day to day. The system consists of one to four trunk routes. Every scout uses one of these trunks in the beginning and end of its food search. In this way, the trunks get continuous chemical reinforcement and do not change much. Each trunk divides into many branch routes. These will change based on food availability.
The organization of foraging trails is strongly affected by the health of the colony and the availability of food. Food deprivation induces a higher amount of foraging ant traffic, compared to a non-deprived population. If a food source is presented to the food deprived colony, this traffic was further increased, an indication of the pharaoh ant's recruitment tactic. If food is not present, a colony will extend its trails to a wider radius around the nest. Logically, number of trails and forager traffic is largest near a food source.
While pheromones explain the pharaoh ant's foraging ability, each forager's capability to find its way back to the nest requires a different explanation. In fact, the pharaoh ant relies on geometry to show it the way home. Each fork in the trail system spreads at an angle between 50 and 60 degrees. When returning to the nest, a forager that encounters a fork will almost always take the path that deviates less from its current direction. In other words, it will never choose an acute angle that would drastically change its direction. Using this algorithm, each forager is able to find its way back to the nest. If the fork angle is experimentally increased to an angle between 60 and 120 degrees, M. pharaonis foragers were significantly less able to find their nest. This method of decision-making reduces the wasted energy that would result from traveling in the wrong direction and contributes to the pharaoh ant's efficiency in foraging.
Feeding
Upon scouts' return with food, the queens will attempt to beg for their share. Depending on food availability and each individual's condition, a scout may refuse the queen's entreaties and even run away from her. The decision of an individual to give up food to the queen may be beneficial in situations of plentiful food, as a healthy queen can reproduce and propagate the colony's genes. However, when food is highly scarce, an individual's own survival can outweigh this potential benefit. She will therefore refuse to give up food.
A queen may also feed on secretions from larvae. This creates a positive feedback loop in which more larvae will provide more food to queens who can in turn produce more larvae.
If a large amount of larvae results in a surplus of secretions, pharaoh ants will store the excess in the gasters of a unique caste, the replete workers. Members of this group have enormous gasters and can regurgitate their stored food when needed. In this way, the colony has a cushion against food shortages.
Pharaoh ants have a sophisticated strategy for food preference. They implement two related behaviors. The first is known as satiation. The workers will at first show a strong preference for a particular food type. However, if this food is offered alone, with no other options, for several weeks, workers will afterward show a distinct preference for a different type of food. In this way, the ants become satiated on a certain food group and will change their decision. The second behavior is called alternation. If given the continuous choice between food groups, pharaoh ants will tend to alternate between carbohydrate-rich foods and protein-rich foods. These satiation and alternation behaviors are evolutionarily adaptive. The decision to vary the type of food consumed ensures that the colony maintains a balanced diet. Edwards & Abraham 1990's result is appropriate for highly competitive environments, and consistent with a high intake:expenditure ratio.
Caste system
Monomorium pharaonis, similar to other invasive ants, is polygynous, meaning its colonies contain many queens (up to 200). It is hypothesized that polygyny leads to lower levels of nestmate recognition in comparison to monogynous species due to the expected higher levels of genetic diversity. Because these colonies lack nestmate recognition, there is no hostility between neighbouring colonies, which is known as unicoloniality.
Many invasive ants display unicoloniality. The adaptive value of this nonaggression among colonies has to do with avoiding unnecessary injury and allowing proper resource allocation, ensuring success for all the colonies. Low nestmate recognition, caused in part by polygyny, also has a biochemical basis in M. pharaonis. Cuticular hydrocarbons are compounds, often found on antennae, that allow for communication in many social insects. In ant species, these compounds play an especially key role in nestmate recognition. Differences in cuticular hydrocarbons are detected by other ant species, who respond accordingly. However, all pharaoh ant colonies have the same hydrocarbons on their antennae. This leads to ineffective nestmate recognition, and nonaggression between colonies.
Pharaoh ant colonies contain many queens. The ratio of queens to workers is variable and dependent on the size of the colony. An individual colony normally contains 1,000–2,500 workers, but often a high density of nests gives the impression of massive colonies. In a small colony, there will be more queens relative to workers. In addition, individuals will be larger than those in a more populous colony. This ratio is controlled by the workers in the colony. Larvae that will produce workers have characteristic hairs all over them, while larvae that will produce sexual males or females are bare. It is thought that workers can use these distinguishing features to identify larvae. Workers may cannibalize larvae in order to ensure a favorable caste ratio. This decision to cannibalize is largely determined by the present caste ratio. If plenty of fertile queens are present, for example, the workers may eat sexual larvae. The caste ratios are controlled in an attempt to maximize the growth of the colony. For example, in a small colony, the ratio of queens to workers is increased. This in turn increases the potential for reproduction, allowing colony growth. Conversely, in a large colony, the high worker to queen ratio maximizes the foraging capacity of the nest, helping sustain the population size.
Nest demographic
The pharaoh ant is a polygynous species that has a relatively low worker to queen ratio of around 12.86. This allows the pharaoh ants to be able to exert social control over the size of the colony and the size of each caste. In the average nest, there are around 170 ± 8 queens, which comprises around 5.2% of the total population, whereas there are around 2185 ± 49 workers, which make up around 66.6% of the population. This low worker to queen ratio is usually associated with swift changes in the nest and may be why pharaoh ants form many new nest buds quickly. To branch out and form a new bud nest, pharaoh ants need a minimum of 469 ± 28 individuals, which explains how they proliferate so quickly.
Reproduction
Mating for pharaoh ants occurs within the nests with males that are usually not from the colony which ensures genetic diversity. The queen can typically produce eggs in batches of 10 to 12 at once, but can lay up to 400 eggs every time she mates. The eggs that are produced take up to 42 days to mature from an egg to an adult. Each queen within the nest lives between 4 and 12 months.
During copulation, sperm is transferred from male to female inside a spermatophore. There are several theories regarding the adaptive value of using a spermatophore. It contains certain chemicals that may inhibit the female's sex drive. Alternatively, it may physically plug the female's gonophore. In either explanation, the spermatophore prevents the female from reproducing with another male. In essence, the use of a spermatophore is evolutionarily favorable because it increases the probability of the male's genetic code being transferred to subsequent generations by lessening potential competition from other males.
Pharaoh ant copulation, like that of many social insects, is harmful to the female. The penis valve contains sharp teeth, which latch onto a thick, soft cuticular layer in the female. This method of copulation too has an evolutionary basis. The teeth ensure sex lasts long enough for sufficient sperm transfer. Also, the pain caused to the female may, in some ways, lessen her desire to mate again.
Queen–worker relationship
When the queen ant first moves to a new nest, she will rear the first group of workers. Once a worker threshold has been reached, resources will then be invested into new males and queens. When a new nest is formed, queens are not a necessity; workers can raise new queens after finding a suitable nest site.
In pharaoh ant colonies new males and queen ants can be produced when the existing fertile queen is removed. When queens are absent, the workers in the nest can do two things: either rear existing sexual larvae or transport sexual larvae from other bud nests or from the main nest to its own nest. However, when there are fertile queens still within the nest, the worker ants will cannibalize the sexual larvae and will either reject or consume sexual larvae from other nests. On the other hand, the worker ants will always accept and nurture worker larvae from other nests. Furthermore, according to Schmidt et al., polygamous species such as pharaoh ants will have higher resource allocations towards the female caste instead of the worker caste to ensure rapid growth of new budding colonies.
Colony interaction
When social ants encounter ants from another colony, behavior can be either aggressive or non-aggressive. Aggressive behavior is very commonly seen; the attacking worker usually bites the opponent at the petiole. In non-aggressive behavior, antennation occurs when the two ants meet. In the case of Monomorium pharaonis, behavior is almost always non-aggressive even when the ants are from different colonies and of different castes. Very few cases exist where aggressive behavior is seen in these ants.
Washing
After foraging, pharaoh ants will wash themselves when a worker enters the cell. Pharaoh ants will also wash after a long feed. It has been proposed that washing has a hygienic value, keeping the nest area clean, staving off disease and disorder. Right before workers leave to forage, they also may wash themselves. However, in this instance the behavior is extremely violent, often causing the ants to fall over. It is thought that here, the washing behavior has no hygienic value and instead may be a displacement activity, a sign that the ants are deliberating whether or not to exit the nest.
Invasiveness and extermination
Budding is a major factor underlying the invasiveness of pharaoh ants. A single seed colony can populate a large office block, almost to the exclusion of all other insect pests, in less than six months. Elimination and control are difficult because multiple colonies can consolidate into smaller colonies during extermination programs only to repopulate later.
Pharaoh ants have become a serious pest in almost every type of building. They can feed on a wide variety of foods including grease, sugary foods, and dead insects. They can also gnaw holes in silk, rayon and rubber goods. Nests can be very small, making detection even more difficult. They are usually found in wall voids, under floors, or in various types of furniture. In homes, they are often found foraging in bathrooms or near food.
It is recommended not to attempt extermination using insecticidal sprays and dusts because they will cause the pharaoh ants to scatter and colonies to split, although non-repellent residual insecticides have been reported to be effective.
The recommended method to eliminate pharaoh ants is by the use of baits attractive to the species. Modern baits use insect growth regulators (IGRs) as the active substance; the ants are attracted to the bait by its food content, and take it back to the nest. Over a period of weeks the IGR prevents the production of worker ants and sterilizes the queen. Renewing the baits once or twice may be necessary.
Pharaoh and other ants have also been exterminated with baits of 1% boric acid and sugar water.
| Biology and health sciences | Hymenoptera | Animals |
41512342 | https://en.wikipedia.org/wiki/USB-C | USB-C | USB-C, or USB Type-C, is a 24-pin connector (not a protocol) that supersedes previous USB connectors and can carry audio, video, and other data, to connect to monitors or external drives. It can also provide and receive power, to power, e.g., a laptop or a mobile phone. It is used not only by USB technology, but also by other protocols, including Thunderbolt, PCIe, HDMI, DisplayPort, and others. It is extensible to support future protocols.
The design for the USB-C connector was initially developed in 2012 by Intel, HP Inc., Microsoft, and the USB Implementers Forum. The Type-C Specification 1.0 was published by the USB Implementers Forum (USB-IF) on August 11, 2014. In July 2016, it was adopted by the IEC as "IEC 62680-1-3".
The USB Type-C connector has 24 pins and is reversible. The designation "C" distinguishes it from the various USB connectors it replaced, all termed either Type-A or Type-B. Whereas earlier USB cables had a host end A and a peripheral device end B, a USB-C cable connects either way; and for interoperation with older equipment, there are cables with a Type-C plug at one end and either a Type-A (host) or a Type-B (peripheral device) plug at the other. The designation "C" refers only to the connector's physical configuration, or form factor, not to be confused with the connector's specific capabilities, such as Thunderbolt 3, DisplayPort 2.0, or USB 3.2 Gen 2x2. Based on the protocols supported by both host and peripheral devices, a USB-C connection normally provides much higher signalling and data rates than the superseded connectors.
A device with a Type-C connector does not necessarily implement any USB transfer protocol, USB Power Delivery, or any of the Alternate Modes: the Type-C connector is common to several technologies while mandating only a few of them.
USB 3.2, released in September 2017, fully replaced the USB 3.1 and USB 3.0 specifications. It preserves the former USB 3.1 SuperSpeed and SuperSpeed+ data transfer modes and introduces two additional data transfer modes by newly applying two-lane operations, with signalling rates of 10 Gbit/s (SuperSpeed USB 10 Gbps; raw data rate: 1.212 GB/s) and 20 Gbit/s (SuperSpeed USB 20 Gbps; raw data rate: 2.422 GB/s). They are only applicable with Full-Featured USB-C Fabrics (connectors, cables, hubs, host, and peripheral device) at all connections.
USB4, released in 2019, is the first USB transfer protocol standard that is applicable exclusively via USB-C.
Ease of use
The USB-C standard simplifies usage by specifying cables having identical plugs on both ends, which can be inserted without concern about orientation. When connecting two devices, the user can plug either end of the cable into either device. The plugs are flat, but will work if inserted right-side-up or upside-down.
The USB-C plugs have two-fold rotational symmetry because a plug may be inserted into a receptacle in either of two orientations. Electrically, USB-C plugs are not symmetric, as can be seen in the tables of pin layouts. Also, the two ends of the USB-C are electrically different, as can be seen in the table of cable wiring. The illusion of symmetry results from how devices respond to the cable. Software makes the plugs and cables behave as though they are symmetric. According to the specifications, "Determination of this host-to-device relationship is accomplished through a Configuration Channel (CC) that is connected through the cable."
The USB-C standard attempts to eliminate the need to have different cables for other communication technologies, such as Thunderbolt, PCIe, HDMI, DisplayPort, Wi-Fi and more. Over the past decade, many companies have adopted the USB-C standard into their products. USB-C cables can contain circuit boards and processors giving them much more capability than simple circuit connections.
Overview
USB-C cables interconnect hosts and peripheral devices, replacing various other electrical cables and connectors, including all earlier (legacy) USB connectors, HDMI connectors, DisplayPort ports, and 3.5 mm audio jacks.
Name
USB Type-C and USB-C are trademarks of the USB Implementers Forum.
Connectors
The 24-pin double-sided connector is slightly larger than the micro-B connector, with a USB-C receptacle measuring wide, high, and deep.
Cables
Type-C cables can be split among various categories and subcategories. The first one is USB 2.0 or Full-Featured. Like the names imply, USB 2.0 Type-C cables have very limited wires and are only good for USB 2.0 communications and power delivery. They are also called charging cables colloquially. Conversely, Full-Featured cables need to have all wires populated and in general support Alt modes and are further distinguished by their speed rating.
Full-Featured cables exist in 4 different speed grades. Their technical names use the "Gen A" notation, each higher number increasing capabilities in terms of bandwidth. The user-facing names are based on the bandwidth a user can typically expect "USB 5Gbps", "USB 20Gbps", "USB 40Gbps" and so on. This bandwidth notation considers the various USB standards and how they use the cable. A Gen 1 / 5 Gbit/s cable supports that bandwidth on every one of its 4 wire pairs. So technically it could be used to establish a USB3 Gen 1x2 connection with nominally 10 Gbit/s between two "SuperSpeed USB 20 Gbps" capable hosts. For a similar reason, the "USB 10Gbps" name is deprecated, as that is using only 2 of the 4 wire-pairs of a Gen 2 cable and thus synonymous with "USB 20Gbps" cables. The signal quality that the "Gen A" notation guarantees or requires is not uniform across all USB standards. See table for details.
The USB Implementers Forum certifies valid cables so they can be marked accordingly with the official logos and users can distinguish them from non-compliant products. There have been simplifications in the logos. Previous logos and names also referenced specific USB protocols like SuperSpeed for the USB3 family of connections or USB4 directly. The current official names and logos have removed those references as most full-featured cables can be used for USB4 connections as well as USB3 connections.
In order to achieve longer cable lengths, cable variants with active electronics to amplify the signals also exist. The Type-C standard mostly mandates these active cables to behave similar to passive cables with vast backwards compatibility. But they are not mandated to support all possible features and typically have no forward compatibility to future standards. Optical cables are even allowed to further reduce the backwards compatibility. For example, an active cable may not be able to use all high speed wire-pairs in the same direction (as used for DisplayPort connections), but only in the symmetric combinations expected by classic USB connections. Passive cables have no such limitations.
Power delivery
Every normal USB-C cable must support at least 3 amps of current and up to 20 volts for up to 60 watts of power according to the USB PD specification. Cables are also allowed to support up to 5 A (with 20 V limit up to 100 W of power). However, the 20 V limit for 5 A cables has been deprecated in favor of 50 V. The combination of higher voltage support and 5 A current support is called and allows for up to 240 W (48 V, 5 A) of power according to the USB PD specification.
E-Marker
All Type-C cables except the minimal combination of USB 2.0 and only 3 A must contain E-Marker chips that identify the cable and its capabilities via the USB PD protocol. This identification data includes information about product/vendor, cable connectors, USB signalling protocol (2.0, Gen speed rating , Gen 2), passive/active construction, use of VCONN power, available VBUS current, latency, RX/TX directionality, SOP controller mode, and hardware/firmware version. It also can include further vendor-defined messages (VDM) that detail support for Alt modes or vendor specific functionality outside of the USB standards.
Cable types
Hosts and peripheral devices
For any two pieces of equipment connecting over USB, one is a host (with a downstream-facing port, DFP) and the other is a peripheral device (with an upstream-facing port, UFP). Some products, such as mobile phones, can take either role, whichever is opposite that of the connected equipment. Such equipment is said to have Dual-Role-Data (DRD) capability, which was known as USB On-The-Go in the previous specification. With USB-C, when two such devices are connected, the roles are first randomly assigned, but a swap can be commanded from either end, although there are optional path and role detection methods that would allow equipment to select a preference for a specific role. Furthermore, Dual-Role equipment that implements USB Power Delivery may swap data and power roles independently using the Data Role Swap or Power Role Swap processes. This allows for charge-through hub or docking station applications such as a portable computer acting as a host to connect to peripherals but being powered by the dock, or a computer being powered by a display, through a single USB-C cable.
USB-C devices may optionally provide or consume bus power currents of 1.5 A and 3.0 A (at 5 V) in addition to baseline bus power provision; power sources can either advertise increased USB current through the configuration channel or implement the full USB Power Delivery specification using both the BMC-coded configuration line and the legacy BFSK-coded VBUS line.
All older USB connectors (all Type-A and Type-B) are designated legacy. Connecting legacy and modern, USB-C equipment requires either a legacy cable assembly (a cable with any Type-A or Type-B plug on one end and a Type-C plug on the other) or, in very specific cases, a legacy adapter assembly.
An older device can connect to a modern (USB-C) host by using a legacy cable, with a Standard-B, Mini-B, or Micro-B plug on the device end and a USB-C plug on the other. Similarly, a modern device can connect to a legacy host by using a legacy cable with a USB-C plug on the device end and a Standard-A plug on the host end. Legacy adapters with USB-C receptacles are "not defined or allowed" by the specification because they can create "many invalid and potentially unsafe" cable combinations (being any cable assembly with two A ends or two B ends). However, exactly three types of adapter with USB-C plugs are defined: 1. A Standard-A receptacle (for connecting a legacy device (such as a flash drive—not a cable) to a modern host, and supporting up to USB 3.1). 2. A Micro-B receptacle (for connecting a modern device to a legacy host, and supporting up to USB 2.0). 3. The Audio adapter accessory mode defined below, in the next section.
Non-USB modes
Audio adapter accessory mode
A device with a USB-C port may support analog headsets through an audio adapter with a 3.5 mm jack, providing three analog audio channels (left and right output and microphone). The audio adapter may optionally include a USB-C charge-through port to allow 500 mA device charging. The engineering specification states that an analog headset shall not use a USB-C plug instead of a 3.5 mm plug. In other words, headsets with a USB-C plug should always support digital audio (and optionally the accessory mode).
Analog signals use the USB 2.0 differential pairs (Dp and Dn for Right and Left) and the two side-band use pairs for Mic and GND. The presence of the audio accessory is signaled through the configuration channel and VCONN.
Alternate modes
An Alternate Mode dedicates some of the physical wires in a USB-C cable for direct device-to-host transmission using non-USB data protocols, such as DisplayPort or Thunderbolt. The four high-speed lanes, two side-band pins, and (for dock, detachable device and permanent-cable applications only) five additional pins can be used for Alternate Mode transmission. The modes are configured using vendor-defined messages (VDM) through the configuration channel.
Specifications
USB Type-C cable and connector specifications
The USB Type-C specification 1.0 was published by the USB Implementers Forum (USB-IF) and was finalized in August 2014.
It defines requirements for cables and connectors.
Rev 1.1 was published 2015-04-03
Rev 1.2 was published 2016-03-25
Rev 1.3 was published 2017-07-14
Rev 1.4 was published 2019-03-29
Rev 2.0 was published 2019-08-29
Rev 2.1 was published 2021-05-25 (USB PD - Extended Power Range - 48 V - 5 A - 240 W)
Rev 2.2 was published 2022-10-18, primarily for enabling USB4 Version 2.0 (80 Gbps) over USB Type-C connectors and cables.
Rev 2.3 was published 2023-10-31.
Adoption as IEC specification:
IEC 62680-1-3:2016 (2016-08-17, edition 1.0) "Universal serial bus interfaces for data and power – Part 1-3: Universal Serial Bus interfaces – Common components – USB Type-C cable and connector specification"
IEC 62680-1-3:2017 (2017-09-25, edition 2.0) "Universal serial bus interfaces for data and power – Part 1-3: Common components – USB Type-C Cable and Connector Specification"
IEC 62680-1-3:2018 (2018-05-24, edition 3.0) "Universal serial bus interfaces for data and power – Part 1-3: Common components – USB Type-C Cable and Connector Specification"
Receptacles
The receptacle features four power and four ground pins, two differential pairs (connected together on devices) for legacy USB 2.0 high-speed data, four shielded differential pairs for Enhanced SuperSpeed data (two transmit and two receive pairs), two Sideband Use (SBU) pins, and two Configuration Channel (CC) pins.
Plugs
The plug has only one USB 2.0 high-speed differential pair, and one of the CC pins (CC2) is replaced by VCONN, to power optional electronics in the cable, and the other is used to actually carry the Configuration Channel (CC) signals. These signals are used to determine the orientation of the cable, as well as to carry USB Power Delivery communications.
Cables
Although plugs have 24 pins, cables have only 18 wires. In the following table, the "No." column shows the wire number.
Related USB-IF specifications
USB Type-C Locking Connector Specification The USB Type-C Locking Connector Specification was published 2016-03-09. It defines the mechanical requirements for USB-C plug connectors and the guidelines for the USB-C receptacle mounting configuration to provide a standardized screw lock mechanism for USB-C connectors and cables.
USB Type-C Port Controller Interface Specification The USB Type-C Port Controller Interface Specification was published 2017-10-01. It defines a common interface from a USB-C Port Manager to a simple USB-C Port Controller.
USB Type-C Authentication Specification Adopted as IEC specification: IEC 62680-1-4:2018 (2018-04-10) "Universal Serial Bus interfaces for data and power - Part 1-4: Common components - USB Type-C Authentication Specification"
USB 2.0 Billboard Device Class Specification USB 2.0 Billboard Device Class is defined to communicate the details of supported Alternate Modes to the computer host OS. It provides user readable strings with product description and user support information. Billboard messages can be used to identify incompatible connections made by users. They optionally appear to negotiate multiple Alternate Modes and must appear when negotiation fails between the host (source) and device (sink).
USB Audio Device Class 3.0 Specification USB Audio Device Class 3.0 defines powered digital audio headsets with a USB-C plug. The standard supports the transfer of both digital and analog audio signals over the USB port.
USB Power Delivery Specification While it is not necessary for USB-C compliant devices to implement USB Power Delivery, for USB-C DRP/DRD (Dual-Role-Power/Data) ports, USB Power Delivery introduces commands for altering a port's power or data role after the roles have been established when a connection is made.
USB 3.2 Specification USB 3.2, released in September 2017, replaces the USB 3.1 specification. It preserves existing USB 3.1 SuperSpeed and SuperSpeed+ data modes and introduces two new SuperSpeed+ transfer modes over the USB-C connector using two-lane operation, doubling the signalling rates to 10 and 20 Gbit/s (raw data rate 1 and ~2.4 GB/s). USB 3.2 is only supported by USB-C, making previously used USB connectors obsolete.
USB4 Specification The USB4 specification released in 2019 is the first USB data transfer specification to be exclusively applicable by the Type-C connector.
Alternate Mode partner specifications
five system-defined Alternate Mode partner specifications exist. Additionally, vendors may support proprietary modes for use in dock solutions. Alternate Modes are optional; Type-C features and devices are not required to support any specific Alternate Mode. The USB Implementers Forum is working with its Alternate Mode partners to make sure that ports are properly labelled with respective logos.
Other protocols like Ethernet have been proposed, although Thunderbolt 3 and later are also capable of 10 Gigabit Ethernet networking.
All Thunderbolt 3 controllers support both Thunderbolt Alternate Mode and DisplayPort Alternate Mode. Because Thunderbolt can encapsulate DisplayPort data, every Thunderbolt controller can either output DisplayPort signals directly over DisplayPort Alternative Mode or encapsulated within Thunderbolt in Thunderbolt Alternate Mode. Low-cost peripherals mostly connect via DisplayPort Alternate Mode while some docking stations tunnel DisplayPort over Thunderbolt.
DisplayPort Alternate Mode 2.0: DisplayPort 2.0 can run directly over USB-C alongside USB4. DisplayPort 2.0 can support 8K resolution at 60 Hz with HDR10 color and can use up to 80 Gbps, which is double the amount available to USB data.
The USB SuperSpeed protocol is similar to DisplayPort and PCIe/Thunderbolt, in using packetized data transmitted over differential LVDS lanes with embedded clock using comparable bit rates, so these Alternate Modes are easier to implement in the chipset.
Alternate Mode hosts and peripheral devices can be connected with either regular Full-Featured Type-C cables, or with converter cables or adapters:
USB 3.1 Type-C to Type-C Full-Featured cable DisplayPort, Mobile High-Definition Link (MHL), HDMI and Thunderbolt (20Gbit/s, or 40Gbit/s with cable length up to 0.5 m) Alternate Mode Type-C ports can be interconnected with standard passive Full-Featured USB Type-C cables. These cables are only marked with standard "trident" SuperSpeed USB logo (for Gen 1 mode only) or the SuperSpeed+ USB 10 Gbit/s logo on both ends. Cable length should be 2.0m or less for Gen 1 and 1.0m or less for Gen 2.
Thunderbolt Type-C to Type-C active cable Thunderbolt 3 (40Gbit/s) Alternate Mode with cables longer than 0.8 m requires active Type-C cables that are certified and electronically marked for high-speed Thunderbolt 3 transmission, similarly to high-power 5 A cables. These cables are marked with a Thunderbolt logo on both ends. They do not support USB 3 backwards compatibility, only USB 2 or Thunderbolt. Cables can be marked for both Thunderbolt and 5 A power delivery at the same time.
Active cables and adapters contain powered electronics to allow for longer cables or to perform protocol conversion. The adapters for video Alternate Modes may allow conversion from native video stream to other video interface standards (e.g., DisplayPort, HDMI, VGA or DVI).
Using Full-Featured Type-C cables for Alternate Mode connections provides some benefits. Alternate Mode does not employ USB 2.0 lanes and the configuration channel lane, so USB 2.0 and USB Power Delivery protocols are always available. In addition, DisplayPort and MHL Alternate Modes can transmit on one, two, or four SuperSpeed lanes, so two of the remaining lanes may be used to simultaneously transmit USB 3.1 data.
USB-C receptacle pin usage in different modes
The diagrams below depict the pins of a USB-C receptacle in different use cases.
USB 2.0/1.1
A simple USB 2.0/1.1 device mates using one pair of D+/D− pins. Hence, the source (host) does not require any connection management circuitry, but it lacks the same physical connector so therefore USB-C is not backward compatible. V and GND provide 5V up to 500mA of current.
However, to connect a USB 2.0/1.1 device to a USB-C host, use of pull-down resistors Rd on the CC pins is required, as the source (host) will not supply V until a connection is detected through the CC pins.
This means many USB-A–to–USB-C cables will only work in the A to C direction (connecting to a USB-C devices, e.g. for charging) as they do not include the termination resistors needed to work in the C to A direction (from a USB-C host). Adapters or cables from USB-C to a USB-A receptacle usually do work as they include the required termination resistor.
USB Power Delivery
The USB Power Delivery specification uses one of CC1 or CC2 pins for power negotiation between source device and sink device, up to 20 V at 5 A. It is transparent to any data transmission mode, and can therefore be used together with any of them as long as the CC pins are intact.
An extension to the specification has added 28 V, 36 V and 48 V to support up to 240 W of power for laptops, monitors, hard disks and other peripherals.
USB 3.0/3.1/3.2
In the USB 3.0/3.1/3.2 mode, two or four high speed links are used in TX/RX pairs to provide 5, 10, or 20 Gbit/s (only by USB 3.2 x2 two-lane operations) signalling rates respectively. One of the CC pins is used to negotiate the mode.
V and GND provide 5 V up to 900 mA, in accordance with the USB 3.1 specification. A specific USB-C mode may also be entered, where 5 V at nominal either 1.5 A or 3 A is provided. A third alternative is to establish a USB Power Delivery (USB-PD) contract.
In single-lane mode, only the differential pairs closest to the CC pin are used for data transmission. For dual-lane data transfers, all four differential pairs are in use.
The D+/D− link for USB 2.0/1.1 is typically not used when a USB 3.x connection is active, but devices like hubs open simultaneous 2.0 and 3.x uplinks in order to allow operation of both types of devices connected to it. Other devices may have the ability to fall back to 2.0, in case the 3.x connection fails. For this, it is important that SS and HS lanes are correctly aligned so that i.e. operating system messages indicating overcurrent conditions report the correct shared USB plug.
Alternate Modes
In Alternate Modes one of up to four high speed links are used in whatever direction is needed. SBU1, SBU2 provide an additional lower speed link. If two high speed links remain unused, then a USB 3.0/3.1 link can be established concurrently to the Alternate Mode. One of the CC pins is used to perform all the negotiation. An additional low band bidirectional channel (other than SBU) may share that CC pin as well. USB 2.0 is also available through D+/D− pins.
In regard to power, the devices are supposed to negotiate a Power Delivery contract before an Alternate Mode is entered.
Debug Accessory Mode
The external device test system (DTS) signals to the target system (TS) to enter debug accessory mode via CC1 and CC2 both being pulled down with an Rd resistor value or pulled up as Rp resistor value from the test plug (Rp and Rd defined in Type-C specification).
After entering debug accessory mode, optional orientation detection via the CC1 and CC2 is done via setting CC1 as a pullup of Rd resistance and CC2 pulled to ground via Ra resistance (from the test system Type-C plug). While optional, orientation detection is required if USB Power Delivery communication is to remain functional.
In this mode, all digital circuits are disconnected from the connector, and the 14 bold pins can be used to expose debug related signals (e.g. JTAG interface). USB IF requires for certification that security and privacy consideration and precaution has been taken and that the user has actually requested that debug test mode be performed.
If a reversible Type-C cable is required but Power Delivery support is not, the test plug will need to be arranged as below, with CC1 and CC2 both being pulled down with an Rd resistor value or pulled up as Rp resistor value from the test plug:
This mirroring of test signals will only provide 7 test signals for debug usage instead of 14, but with the benefit of minimizing extra parts count for orientation detection.
Audio Adapter Accessory Mode
In this mode, all digital circuits are disconnected from the connector, and certain pins become reassigned for analog outputs or inputs. The mode, if supported, is entered when both CC pins are shorted to GND. D− and D+ become audio output left L and right R, respectively. The SBU pins become a microphone pin MIC, and the analog ground AGND, the latter being a return path for both outputs and the microphone. Nevertheless, the MIC and AGND pins must have automatic swap capability, for two reasons: firstly, the USB-C plug may be inserted either side; secondly, there is no agreement, which TRRS rings shall be GND and MIC, so devices equipped with a headphone jack with microphone input must be able to perform this swap anyway.
This mode also allows concurrent charging of a device exposing the analog audio interface (through V and GND), however only at 5 V and 500 mA, as CC pins are unavailable for any negotiation.
Plug insertions detection is performed by the TRRS plug's physical plug detection switch. On plug insertions, this will pull down both CC and VCONN in the plug (CC1 and CC2 in the receptacle). This resistance must be less than 800 ohms which is the minimum "Ra" resistance specified in the USB Type-C specification). This is essentially a direct connection to USB digital ground.
Software support
Android from version 6.0 "Marshmallow" onwards works with USB 3.1 and USB-C.
ChromeOS, starting with the Chromebook Pixel 2015, supports USB 3.1, USB-C, Alternate Modes, Power Delivery, and USB Dual-Role support.
FreeBSD released the Extensible Host Controller Interface, supporting USB 3.0, with release 8.2
iOS from version 12.1 (iPad Pro 3rd generation or later, iPad Air 4th generation or later, iPad Mini 6th generation or later, iPad 10th generation or later, iPhone 15 or later) onwards works with USB-C.
NetBSD began supporting USB 3.0 with release 7.2
Linux has supported USB 3.0 since kernel version 2.6.31 and USB version 3.1 since kernel version 4.6.
OpenBSD began supporting USB 3.0 in version 5.7
OS X Yosemite (macOS version 10.10.2), starting with the MacBook Retina early 2015, supports USB 3.1, USB-C, Alternate Modes, and Power Delivery.
Windows 8.1 added USB-C and billboard support in an update.
Windows 10 and Windows 10 Mobile support USB 3.1, USB-C, alternate modes, billboard device class, Power Delivery and USB Dual-Role.
Authentication
USB Type-C Authentication is an extension to the USB-C protocol which can add security to the protocol.
Hardware support
USB-C devices
An increasing number of motherboards, notebooks, tablet computers, smartphones, hard disk drives, USB hubs and other devices released from 2014 onwards include the USB-C receptacles. However, the initial adoption of USB-C was limited by the high cost of USB-C cables and the wide use of Micro-USB chargers.
Video output
Currently, DisplayPort is the most widely implemented alternate mode, and is used to provide video output on devices that do not have standard-size DisplayPort or HDMI ports, such as smartphones and laptops. All Chromebooks with a USB-C port are required to support DisplayPort alternate mode in Google's hardware requirements for manufacturers. A USB-C multiport adapter converts the device's native video stream to DisplayPort/HDMI/VGA, allowing it to be displayed on an external display, such as a television set or computer monitor.
It is also used on USB-C docks designed to connect a device to a power source, external display, USB hub, and optional extra (such as a network port) with a single cable. These functions are sometimes implemented directly into the display instead of a separate dock, meaning a user connects their device to the display via USB-C with no other connections required.
Compatibility issues
Power issues with cables
Many cables claiming to support USB-C are actually not compliant to the standard. These cables can, potentially, damage a device. There are reported cases of laptops being destroyed due to the use of non-compliant cables.
Some non-compliant cables with a USB-C connector on one end and a legacy USB-A plug or Micro-B receptacle (receptacles also usually being invalid on cables, but see known exceptions in the sections on Hosts and peripheral devices and Audio adapter accessory mode above) on the other end incorrectly terminate the Configuration Channel (CC) with a 10 kΩ pull-up to VBUS instead of the specification mandated 56 kΩ pull-up, causing a device connected to the cable to incorrectly determine the amount of power it is permitted to draw from the cable. Cables with this issue may not work properly with certain products, including Apple and Google products, and may even damage power sources such as chargers, hubs, or PC USB ports.
A defective USB-C cable or power source can cause a USB-C device to see and an incorrect and different "declared" voltage than what the source will actually deliver. This may result in an overvoltage on the VBUS pin.
Also due to the fine pitch of the USB-C receptacle, the VBUS pin from the cable may contact with the CC pin of the USB-C receptacle resulting in a short-to-VBUS electrical issue due to the fact that the VBUS pin is rated up to 20 V while the CC pins are rated up to 5.5 V.
To overcome these issues, USB Type-C port protection must be used between a USB-C connector and a USB-C Power Delivery controller.
Compatibility with audio adapters
The USB-C port can be used to connect wired accessories such as headphones.
There are two modes of audio output from devices: digital and analog. There are primarily two types of USB-C audio adapters: active, e.g. those with digital-to-analog converters (DACs), and passive, without electronics.
When an active set of USB-C headphones or adapter is used, digital audio is sent through the USB-C port. The conversion by the DAC and amplifier is done inside of the headphones or adapter, instead of on the phone. The sound quality is dependent on the headphones/adapter's DAC. Active adapters with a built-in DAC have near-universal support for devices that output digital and analog audio, adhering to the Audio Device Class 3.0 and Audio Adapter Accessory Mode specifications.
Examples of such active adapters include external USB sound cards and DACs that do not require special drivers, and USB-C to 3.5 mm headphone jack adapters by Apple, Google, Essential, Razer, HTC, and Samsung.
On the other hand, when a passive adapter is used, digital-to-analog conversion is done on the host device and analog audio is sent through the USB-C port. The sound quality is dependent on the phone's onboard DAC. Passive adapters are only compatible with devices that output analog audio, adhering to the Audio Adapter Accessory Mode specification.
Compatibility with other fast-charging technology
In 2016, Benson Leung, an engineer at Google, pointed out that Quick Charge 2.0 and 3.0 technologies developed by Qualcomm are not compatible with the USB-C standard. Qualcomm responded that it is possible to make fast-charge solutions fit the voltage demands of USB-C and that there are no reports of problems; however, it did not address the standard compliance issue at that time. Later in the year, Qualcomm released Quick Charge 4, which it claimed was – as an advancement over previous generations – "USB Type-C and USB PD compliant".
Regulations for compatibility
In 2021, the European Commission proposed the use of USB-C as a universal charger. On 4 October 2022, the European Parliament voted in favor of the new law, Radio Equipment Directive 2022/2380, with 602 votes in favor, 13 against and 8 abstentions. The regulation requires that all new mobile phones, tablets, cameras, headphones, headsets, handheld video game consoles, portable speakers, e-readers, keyboards, mice, portable navigation systems, and earbuds sold in the European Union and supporting wired charging, would have to be equipped with a USB-C port and charge with a standard USB-C to USB-C cable by the end of 2024. Additionally, if these devices support fast charging, they must support USB Power Delivery. These regulations will extend to laptops by early 2026. To comply with these regulations, Apple Inc. replaced its proprietary Lightning connector with USB-C beginning with the iPhone 15 and AirPods Pro second generation, released in 2023. A first modified iPhone having USB-C connector was the result of a hack by Ken Pillonel.
In late December 2024, new EU regulations took effect, mandating USB-C charging ports for all small and medium-sized electronic devices sold in the EU, with laptops to follow by 2026. These rules were aimed at reducing waste and saving €250 million annually for consumers. Apple, which initially opposed the changes, had since adopted USB-C for its products. Additionally, consumers can opt not to receive a new charger with their device.
| Technology | User interface | null |
1461205 | https://en.wikipedia.org/wiki/Argon%E2%80%93argon%20dating | Argon–argon dating | Argon–argon (or 40Ar/39Ar) dating is a radiometric dating method invented to supersede potassiumargon (K/Ar) dating in accuracy. The older method required splitting samples into two for separate potassium and argon measurements, while the newer method requires only one rock fragment or mineral grain and uses a single measurement of argon isotopes. 40Ar/39Ar dating relies on neutron irradiation from a nuclear reactor to convert a stable form of potassium (39K) into the radioactive 39Ar. As long as a standard of known age is co-irradiated with unknown samples, it is possible to use a single measurement of argon isotopes to calculate the 40K/40Ar* ratio, and thus to calculate the age of the unknown sample. 40Ar* refers to the radiogenic 40Ar, i.e. the 40Ar produced from radioactive decay of 40K. 40Ar* does not include atmospheric argon adsorbed to the surface or inherited through diffusion and its calculated value is derived from measuring the 36Ar (which is assumed to be of atmospheric origin) and assuming that 40Ar is found in a constant ratio to 36Ar in atmospheric gases.
Method
The sample is generally crushed and single crystals of a mineral or fragments of rock are hand-selected for analysis. These are then irradiated to produce 39Ar from 39K via the (n-p) reaction 39K(n,p)39Ar. The sample is then degassed in a high-vacuum mass spectrometer via a laser or resistance furnace. Heating causes the crystal structure of the mineral (or minerals) to degrade, and, as the sample melts, trapped gases are released. The gas may include atmospheric gases, such as carbon dioxide, water, nitrogen, and radiogenic gases like argon and helium, generated from regular radioactive decay over geologic time. The abundance of 40Ar* increases with the age of the sample, though the rate of increase decays exponentially with the half-life of 40K, which is 1.248 billion years.
Age equation
The age of a sample is given by the age equation:
where λ is the radioactive decay constant of 40K (approximately 5.5 x 10−10 year−1, corresponding to a half-life of approximately 1.25 billion years), J is the J-factor (parameter associated with the irradiation process), and R is the 40Ar*/39Ar ratio. The J factor relates to the fluence of the neutron bombardment during the irradiation process; a denser flow of neutron particles will convert more atoms of 39K to 39Ar than a less dense one.
Relative dating only
The 40Ar/39Ar method only measures relative dates. In order for an age to be calculated by the 40Ar/39Ar technique, the J parameter must be determined by irradiating the unknown sample along with a sample of known age for a standard. Because this (primary) standard ultimately cannot be determined by 40Ar/39Ar, it must be first determined by another dating method. The method most commonly used to date the primary standard is the conventional K/Ar technique. An alternative method of calibrating the used standard is astronomical tuning (also known as orbital tuning), which arrives at a slightly different age.
Applications
The primary use for 40Ar/39Ar geochronology is dating metamorphic and igneous minerals. 40Ar/39Ar is unlikely to provide the age of intrusions of granite as the age typically reflects the time when a mineral cooled through its closure temperature. However, in a metamorphic rock that has not exceeded its closure temperature the age likely dates the crystallization of the mineral. Dating of movement on fault systems is also possible with the 40Ar/39Ar method. Different minerals have different closure temperatures; biotite is ~300°C, muscovite is about 400°C and hornblende has a closure temperature of ~550°C. Thus, a granite containing all three minerals will record three different "ages" of emplacement as it cools down through these closure temperatures. Thus, although a crystallization age is not recorded, the information is still useful in constructing the thermal history of the rock.
Dating minerals may provide age information on a rock, but assumptions must be made. Minerals usually only record the last time they cooled down below the closure temperature, and this may not represent all of the events which the rock has undergone, and may not match the age of intrusion. Thus, discretion and interpretation of age dating is essential. 40Ar/39Ar geochronology assumes that a rock retains all of its 40Ar after cooling past the closing temperature and that this was properly sampled during analysis.
This technique allows the errors involved in K-Ar dating to be checked. Argon–argon dating has the advantage of not requiring determinations of potassium. Modern methods of analysis allow individual regions of crystals to be investigated. This method is important as it allows crystals forming and cooling during different events to be identified.
Recalibration
One problem with argon-argon dating has been a slight discrepancy with other methods of dating. Work by Kuiper et al. reports that a correction of 0.65% is needed. Thus the Cretaceous–Paleogene extinction (when the dinosaurs died out)—previously dated at 65.0 or 65.5 million years ago—is more accurately dated to 66.0-66.1 Ma.
| Physical sciences | Geochronology | Earth science |
1461372 | https://en.wikipedia.org/wiki/Laser%20rangefinder | Laser rangefinder | A laser rangefinder, also known as a laser telemeter, is a rangefinder that uses a laser beam to determine the distance to an object. The most common form of laser rangefinder operates on the time of flight principle by sending a laser pulse in a narrow beam towards the object and measuring the time taken by the pulse to be reflected off the target and returned to the sender. Due to the high speed of light, this technique is not appropriate for high precision sub-millimeter measurements, where triangulation and other techniques are often used instead. Laser rangefinders are sometimes classified as type of handheld scannerless lidar.
Pulse
The pulse may be coded to reduce the chance that the rangefinder can be jammed. It is possible to use Doppler effect techniques to judge whether the object is moving towards or away from the rangefinder, and if so, how fast.
Precision
The precision of an instrument is correlated with the rise time, divergence, and power of its laser pulse, as well as the quality of its optics and onboard digital signal processing. Environmental factors can significantly reduce range and accuracy:
Humidity, snow, dust, or other airborne particulates will diffuse the signal.
Higher temperature and higher pressure (lower elevation) slightly decrease the speed of light through air.
Smaller and less reflective targets return less information.
In good conditions, skilled operators using precision laser rangefinders can range a target to within a meter at distances on the order of three kilometers.
Range and range error
Despite the beam being narrow, it will eventually spread over long distances due to the divergence of the laser beam, as well as due to scintillation and beam wander effects, caused by the presence of water droplets in the air acting as lenses ranging in size from microscopic to roughly half the height of the laser beam's path above the earth.
These atmospheric distortions coupled with the divergence of the laser itself and with transverse winds that serve to push the atmospheric heat bubbles laterally may combine to make it difficult to get an accurate reading of the distance of an object, say, beneath some trees or behind bushes, or even over long distances of more than 1 km in open and unobscured desert terrain.
Some of the laser light might reflect off leaves or branches which are closer than the object, giving an early return and a reading which is too low. Alternatively, over distances longer than 360 m, if the target is in proximity to the earth, it may simply vanish into a mirage, caused by temperature gradients in the air in proximity to the heated surface bending the laser light. All these effects must be considered.
Calculation
The distance between point A and B is given by
where c is the speed of light and t is the amount of time for the round-trip between A and B.
where φ is the phase delay made by the light traveling and ω is the angular frequency of optical wave.
Then substituting the values in the equation,
In this equation, λ is the wavelength ; Δφ is the part of the phase delay that does not fulfill (that is, φ modulo ); N is the integer number of wave half-cycles of the round-trip and ΔN the remaining fractional part.
Technologies
Time of flight - this measures the time taken for a light pulse to travel to the target and back. With the speed of light known, and an accurate measurement of the time taken, the distance can be calculated. Many pulses are fired sequentially and the average response is most commonly used. This technique requires very accurate sub-nanosecond timing circuitry.
Multiple frequency phase-shift - this measures the phase shift of multiple frequencies on reflection then solves some simultaneous equations to give a final measure.
Interferometry - the most accurate and most useful technique for measuring changes in distance rather than absolute distances.
Light attenuation by atmospheric absorption - The method measures the attenuation of a laser beam caused by the absorption from an atmospheric compound (H2O, CO2, CH4, O2 etc.) to calculate the distance to an object. The light atmospheric absorption attenuation method requires unmodulated incoherent light sources and low-frequency electronics that reduce the complexity of the devices. Due to this, low-cost light sources can be used for range-finding. However, the application of the method is limited to atmospheric measurements or planetary exploration.
Applications
Military
Rangefinders provide an exact distance to targets located beyond the distance of point-blank shooting to snipers and artillery. They can also be used for military reconnaissance and engineering. Usually tanks use LRF to correct the direct shoot solution.
Handheld military rangefinders operate at ranges of 2 km up to 25 km and are combined with binoculars or monoculars. When the rangefinder is equipped with a digital magnetic compass (DMC) and inclinometer it is capable of providing magnetic azimuth, inclination, and height (length) of targets. Some rangefinders can also measure a target's speed in relation to the observer. Some rangefinders have cable or wireless interfaces to enable them to transfer their measurement(s) data to other equipment like fire control computers. Some models also offer the possibility to use add-on night vision modules. Most handheld rangefinders use standard or rechargeable batteries.
The more powerful models of rangefinders measure distance up to 40 km and are normally installed either on a tripod or directly on a vehicle, ship, jet, helicopter or gun platform. In the latter case the rangefinder module is integrated with on-board thermal, night vision and daytime observation equipment. The most advanced military rangefinders can be integrated with computers.
To make laser rangefinders and laser-guided weapons less useful against military targets, various military arms may have developed laser-absorbing paint for their vehicles. Regardless, some objects don't reflect laser light very well and using a laser rangefinder on them is difficult.
The first commercial laser rangefinder was the Barr & Stroud LF1, developed in association with Hughes Aircraft, which became available in 1965. This was then followed by the Barr & Stroud LF2, which integrated the rangefinder into a tank sight, and this was used on the Chieftain tank in 1969, the first vehicle so-equipped with such a system. Both systems used ruby lasers.
3D modelling
Laser rangefinders are used extensively in 3D object recognition, 3D object modelling, and a wide variety of computer vision-related fields. This technology constitutes the heart of the so-called time-of-flight 3D scanners. In contrast to the military instruments, laser rangefinders offer high-precision scanning abilities, with either single-face or 360-degree scanning modes.
A number of algorithms have been developed to merge the range data retrieved from multiple angles of a single object to produce complete 3D models with as little error as possible. One of the advantages offered by laser rangefinders over other methods of computer vision is in not needing to correlate features from two images in order to determine depth-information like stereoscopic methods do.
Laser rangefinders used in computer vision applications often have depth resolutions of 0.1 mm or less. This can be achieved by using triangulation or refraction measurement techniques unlike to the time of flight techniques used in LIDAR.
Forestry
Special laser rangefinders are used in forestry. These devices have anti-leaf filters and work with reflectors. Laser beam reflects only from this reflector and so exact distance measurement is guaranteed. Laser rangefinders with anti-leaf filter are used for example for forest inventories.
Sports
Laser rangefinders may be effectively used in various sports that require precision distance measurement, such as golf, hunting, and archery. Some of the more popular manufacturers are Caddytalk, Opti-logic Corporation, Bushnell, Leupold, LaserTechnology, Trimble, Leica, Newcon Optik, Op. Electronics, Nikon, Swarovski Optik and Zeiss.
Many rangefinders from Bushnell come with advanced features, such as ARC (angle range compensation), multi-distance ability, slope, JOLT (Vibrate when the target is locked), and Pin-Seeking. ARC can be calculated by hand using the rifleman's rule, but it's usually much easier if you let a rangefinder do it when you are out hunting. In golfing where time is most important, a laser rangefinder comes useful in locating distance to the flag. However not all features are 100% legal for golf tournament play. Many hunters in the eastern U.S. don't need a rangefinder, although many western hunters need them, due to longer shooting distances and more open spaces.
Industrial production processes
An important application is the use of laser rangefinder technology during the automation of stock management systems and production processes in steel industry.
Laser measuring tools
Laser rangefinders are also used in several industries like construction, renovation and real estate as alternatives to tape measures, and was first introduced by Leica Geosystems in 1993 in France. To measure a large object like a room with a tape measure, one would need another person to hold the tape at the far wall and a clear line straight across the room to stretch the tape. With a laser measuring tool, the job can be completed by one operator with just a line of sight. Although tape measures are technically perfectly accurate, laser measuring tools are much more precise. Laser measuring tools typically include the ability to produce some simple calculations, such as the area or volume of a room. These devices can be found in hardware stores and online marketplaces.
Price
Laser rangefinders can vary in price, depending on the quality and application of the product. Military grade rangefinders need to be as accurate as possible and must also reach great distances. These devices can cost hundreds of thousands of dollars. For civilian applications, such as hunting or golf, devices are more affordable and much more readily accessible.
Safety
Laser rangefinders are divided into four classes and several subclasses. Laser rangefinders available to consumers are usually laser class 1 or class 2 devices and are considered relatively eye-safe. Regardless of the safety rating, direct eye contact should always be avoided. Most laser rangefinders for military use exceed the laser class 2 energy levels.
| Technology | Surveying tools | null |
1463084 | https://en.wikipedia.org/wiki/Persian%20gardens | Persian gardens | The tradition and style of garden design represented by Persian gardens or Iranian gardens () is a style of "landscape" garden which emerged in the Achaemenid Empire. Humayun's Tomb and the Taj Mahal have some of the largest Persian gardens in the world, from the era of the Mughal Empire in India.
Concept and etymology
From the time of the Achaemenid Empire, the idea of an earthly paradise spread through Persian literature and example to other cultures, both the Hellenistic gardens of the Seleucid Empire and the Ptolemies in Alexandria.
The Avestan word pairidaēza-, Old Persian *paridaida-, or Median *paridaiza- "walled-around", (i.e., a walled garden), were borrowed into Elamite (partetaš) and Akkadian, and later as . It was rendered as Latin paradīsus, and from there entered into European languages, e.g., French paradis, German Paradies, and English paradise. In the Achaemenid Empire, the term is used for functional worksites as well, and it translated as "plantations"; these sites contained not only orchards and tree plantations, but also sometimes bitumen harvesting and mining. The same word was used to describe the Elamite custom of creating a sacred grove or husa surrounding a royal grave that was the site of worship of the deceased king.
As the word expresses, such gardens would have been enclosed. Gardens provided a place for protected relaxation, both spiritual and leisurely, and represented a paradise on earth. The Common Iranian word for "enclosed space" was *pari-daiza- (Avestan pairi-daēza-). This term was adopted to describe the garden of Eden a Paradise on earth.
The garden's construction may be formal (with an emphasis on structure) or casual (with a focus on nature), following several simple design rules. This allows the maximum possible use of the garden in terms of function and emotion.
History
Persian gardens may originate as early as 4000 BC, but it is clear that this tradition began with the Achaemenid dynasty around the 6th century BCE. Decorated pottery of that time displays the typical cross plan of the Persian garden. The outline of Pasargadae, built around 500 BC, is still viewable today. Classical Iranians were seen by the Greeks as the 'great gardeners' of antiquity; Cyrus II (known also as Cyrus the Younger) is alleged to have told the Spartan commander Lysander that he gardened daily when not campaigning, and had himself laid out the park at Sardis, which he called his 'paradise' (a Greek corruption of the Old Persian word for garden).
During the suzerainty of the Sasanian Empire, under the influence of Zoroastrianism, water in art grew increasingly important. This trend manifested itself in garden design, with greater emphasis on fountains and ponds in gardens.
During the Umayyad and Abbasid periods, the aesthetic aspect of the garden increased in importance, overtaking utility. During this time, aesthetic rules that govern the garden grew in importance. An example of this is the chahār bāgh (), a form of garden that attempts to emulate the Abrahamic notion of a Garden of Eden, with four rivers and four quadrants that represent the world. The design sometimes extends one axis longer than the cross-axis and may feature water channels that run through each of the four gardens and connect to a central pool.
Under the Abbasid dynasty (8th century AD), this type of garden became an integral part of representational architecture.
The Persian garden is a landscape garden, designed individually and created intentionally as a space embedded in the aesthetic and spiritual context of its past and contemporary cultural, political, and social environment. Hallmarks of these formal gardens are a geometric layout following geometric and visual principles, implemented to nature by water channels and basins which divide the enclosed space into clearly defined quarters, a principle that has become known as Chahar Bagh (four gardens), waterworks with channels, basins, fountains and cascades, pavilions, prominent central axes with a vista, and a plantation with a variety of carefully chosen trees, herbs. and flowers. The old-Iranian word for such gardens "pari-daizi' expresses the notion of an earthly paradise that is inherent to them. As such, they are a metaphor for the divine order and the unification and protection of the ones who do good. Their counterparts on earth fulfill a similar function. These principles are brought to perfection in the gardens of the emperor as the "good gardener".
Notwithstanding a formal standardization, the landscape gardens also reflect diversity and development, bound to function, regional and chronological characteristics, as well as technological, know how personal preferences, ambitions, and demands. Persian gardens are multi-functional: they not only serve contemplation and relaxation, but are also a representation and manifestation of power. Designing and implementing a garden demonstrates the occupation of land, holding audiences and celebrating victories or marriages in these gardens signal superiority, or social and political bonds. Starting from the 12th to 13th century, tombs for members of the royal family or important personalities were placed into such formal gardens, providing believers a chance to benefit from the spirituality of a venerated person and the particular aura of the garden.
The invasion of Persia by the Mongols in the thirteenth century led to a new emphasis on highly ornate structure in the garden. Examples of this include tree peonies and chrysanthemums. The Mongols then carried a Persian garden tradition to other parts of their empire (notably India).
The Mughal emperor Babur introduced the Persian garden to India, attempting to replicate the cool, refreshing aura of his homeland in the Ferghana Valley through the construction of Persian-style gardens, like those at other Timurid cities like Samarkand and Herat. Babur was a zealous gardener and personally designed and supervised at least ten gardens in his capital of Kabul in modern Afghanistan, such as the Bagh-e Babur, where he recorded the allure of the pomegranate, cherry and orange trees he had planted. Though his empire soon expanded as far as north-central India, he abhorred the stagnant heat and drab environment of the hot, dusty plains of India; he was thus interred at Bagh-e Babur in Kabul by his widow in 1544.
The Aram Bagh of Agra was the first of many Persian gardens he created in India itself. Mughal gardens have four basic requirements, symbolizing four allegorical essentials for the afterlife: shade, fruit, fragrance, and running water. This pattern was used to build many Persian gardens throughout the Indian subcontinent, such as the Shalimar Gardens of Lahore, the Shalimar Bagh and Nishat Bagh of Kashmir, and the Taj Mahal gardens. The Taj Mahal gardens embody the Persian concept of an ideal paradise garden, and were built with irrigation channels and canals from the Yamuna River. These gardens have recently been restored to their former beauty after decades of pollution by the Indian authorities, who cut down the fruit- and shade-bearing vegetation of the garden.
The Safavid dynasty (seventeenth to eighteenth century) built and developed grand and epic layouts that went beyond a simple extension to a palace and became an integral aesthetic and functional part of it. In the following centuries, European garden design began to influence Persia, particularly the designs of France, and secondarily that of Russia and the United Kingdom. Western influences led to changes in the use of water and the species used in bedding.
Traditional forms and style are still applied in modern Iranian gardens. They also appear in historic sites, museums, and affixed to the houses of the rich.
Elements of the Persian garden
Sunlight and its effects were an important factor of structural design in Persian gardens. Textures and shapes were specifically chosen by architects to harness the light.
Iran's dry heat makes shade important in gardens, which would be nearly unusable without it. Trees and trellises largely feature as biotic shade; pavilions and walls are also structurally prominent in blocking the sun.
The heat also makes water important, both in the design and maintenance of the garden. Irrigation may be required, and may be provided via a form of tunnel called a qanat, that transports water from a local aquifer. Well-like structures then connect to the qanat, enabling the drawing of water. Alternatively, an animal-driven Persian well would draw water to the surface. Such wheel systems also moved water around surface water systems, such as those in the chahar bāgh style. Trees were often planted in a ditch called a juy, which prevented water evaporation and allowed the water quick access to the tree roots.
The Persian style often attempts to integrate indoors with outdoors through the connection of a surrounding garden with an inner courtyard. Designers often place architectural elements such as vaulted arches between the outer and interior areas to open up the divide between them.
Descriptions
An early description (from the first half of the fourth century BCE) of a Persian garden is found in Xenophon's Oeconomicus in which he has Socrates relate the story of the Spartan general Lysander's visit to the Persian prince Cyrus the Younger, who shows the Greek his "paradise at Sardis". In this story Lysander is "astonished at the beauty of the trees within, all planted at equal intervals, the long straight rows of waving branches, the perfect regularity, the rectangular symmetry of the whole, and the many sweet scents which hung about them as they paced the park"
The oldest representational descriptions and illustrations of Persian gardens come from travelers who reached Iran from the west. These accounts include Ibn Battuta in the fourteenth century, Ruy González de Clavijo in the fifteenth century and Engelbert Kaempfer in the seventeenth century. Battuta and Clavijo made only passing references to gardens and did not describe their design, but Kaempfer made careful drawings and converted them into detailed engravings after his return to Europe. They show charbagh-type gardens that featured an enclosing wall, rectangular pools, an internal network of canals, garden pavilions and lush planting. There are surviving examples of this garden type at Yazd (Dowlatabad) and at Kashan (Fin Garden). The location of the gardens Kaempfer illustrated in Isfahan can be identified.
Styles
The six primary styles of the Persian garden may be seen in the following table, which puts them in the context of their function and style. Gardens are not limited to a particular style, but often integrate different styles, or have areas with different functions and styles.
Hayāt
Publicly, it is a classical Persian layout with heavy emphasis on aesthetics over function. Man-made structures in the garden are particularly important, with arches and pools (which may be used to bathe). The ground is often covered in gravel flagged with stone. Plantings are typically very simple - such as a line of trees, which also provide shade.
Privately, these gardens are often pool-centred and, again, structural. The pool serves as a focus and source of humidity for the surrounding atmosphere. There are few plants, often due to the limited water available in urban areas.
Meidān
This is a public, formal garden that puts more emphasis on the biotic element than the hayāt and that minimises structure. Plants range from trees, to shrubs, to bedding plants, to grasses. Again, there are elements such as a pool and gravel pathways which divide the lawn. When structures are used, they are often built, as in the case of pavilions, to provide shade.
Chahar Bāgh
These gardens are private and formal. The basic structure consists of four quadrants divided by waterways or pathways. Traditionally, the rich used such gardens in work-related functions (such as entertaining ambassadors). These gardens balance structure with greenery, with the plants often around the periphery of a pool and path based structure.
Boostān
Much like many other parks, the Persian park serves a casual public function with emphasis on plant life. They provide pathways and seating, but are otherwise usually limited in terms of structural elements. The purpose of such places is relaxation and socialisation.
Bāgh
Like the other casual garden, the park, bāgh emphasizes the natural and green aspect of the garden. Unlike the park it is a private area often affixed to houses and often consisting of lawns, trees, and ground plants. The waterways and pathways stand out less than in the more formal counterparts and are largely functional. The primary function of such areas is familial relaxation.
World Heritage Sites
Pasargad Garden at Pasargadae, Iran (WHS 1372-001)
Eram Garden, Shiraz, Iran (WHS 1372-002)
Chehel Sotoun, Isfahan, Iran (WHS 1372-003)
Fin Garden, Kashan, Iran (WHS 1372-004)
Abbasabad Garden, Abbasabad, Mazandaran, Iran (WHS 1372-005)
Shazdeh Garden, Mahan, Kerman Province, Iran (WHS 1372-006)
Dolatabad Garden, Yazd, Iran (WHS 1372-007)
Pahlevanpour Garden, Iran (WHS 1372-008)
Akbarieh Garden, South Khorasan Province, Iran (WHS 1372-009)
Taj Mahal, Agra, India (WHS 252)
Humayun's Tomb, New Delhi, India (WHS 232bis)
Shalimar Gardens, Lahore, Pakistan (WHS 171-002)
| Technology | Buildings and infrastructure | null |
1463286 | https://en.wikipedia.org/wiki/Enantioselective%20synthesis | Enantioselective synthesis | Enantioselective synthesis, also called asymmetric synthesis, is a form of chemical synthesis. It is defined by IUPAC as "a chemical reaction (or reaction sequence) in which one or more new elements of chirality are formed in a substrate molecule and which produces the stereoisomeric (enantiomeric or diastereomeric) products in unequal amounts."
Put more simply: it is the synthesis of a compound by a method that favors the formation of a specific enantiomer or diastereomer. Enantiomers are stereoisomers that have opposite configurations at every chiral center. Diastereomers are stereoisomers that differ at one or more chiral centers.
Enantioselective synthesis is a key process in modern chemistry and is particularly important in the field of pharmaceuticals, as the different enantiomers or diastereomers of a molecule often have different biological activity.
Overview
Many of the building blocks of biological systems such as sugars and amino acids are produced exclusively as one enantiomer. As a result, living systems possess a high degree of chemical chirality and will often react differently with the various enantiomers of a given compound. Examples of this selectivity include:
Flavour: the artificial sweetener aspartame has two enantiomers. L-aspartame tastes sweet whereas D-aspartame is tasteless.
Odor: R-(–)-carvone smells like spearmint whereas S-(+)-carvone smells like caraway.
Drug effectiveness: the antidepressant drug Citalopram is sold as a racemic mixture. However, studies have shown that only the (S)-(+) enantiomer is responsible for the drug's beneficial effects.
Drug safety: D‑penicillamine is used in chelation therapy and for the treatment of rheumatoid arthritis whereas L‑penicillamine is toxic as it inhibits the action of pyridoxine, an essential B vitamin.
As such enantioselective synthesis is of great importance but it can also be difficult to achieve. Enantiomers possess identical enthalpies and entropies and hence should be produced in equal amounts by an undirected process – leading to a racemic mixture. Enantioselective synthesis can be achieved by using a chiral feature that favors the formation of one enantiomer over another through interactions at the transition state. This biasing is known as asymmetric induction and can involve chiral features in the substrate, reagent, catalyst, or environment and works by making the activation energy required to form one enantiomer lower than that of the opposing enantiomer.
Enantioselectivity is usually determined by the relative rates of an enantiodifferentiating step—the point at which one reactant can become either of two enantiomeric products. The rate constant, k, for a reaction is function of the activation energy of the reaction, sometimes called the energy barrier, and is temperature-dependent. Using the Gibbs free energy of the energy barrier, ΔG*, means that the relative rates for opposing stereochemical outcomes at a given temperature, T, is:
This temperature dependence means the rate difference, and therefore the enantioselectivity, is greater at lower temperatures. As a result, even small energy-barrier differences can lead to a noticeable effect.
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Approaches
Enantioselective catalysis
Enantioselective catalysis (known traditionally as "asymmetric catalysis") is performed using chiral catalysts, which are typically chiral coordination complexes. Catalysis is effective for a broader range of transformations than any other method of enantioselective synthesis. The chiral metal catalysts are almost invariably rendered chiral by using chiral ligands, but it is possible to generate chiral-at-metal complexes composed entirely of achiral ligands.
Most enantioselective catalysts are effective at low substrate/catalyst ratios. Given their high efficiencies, they are often suitable for industrial scale synthesis, even with expensive catalysts. A versatile example of enantioselective synthesis is asymmetric hydrogenation, which is used to reduce a wide variety of functional groups.
The design of new catalysts is dominated by the development of new classes of ligands. Certain ligands, often referred to as "privileged ligands", are effective in a wide range of reactions; examples include BINOL, Salen, and BOX. Most catalysts are effective for only one type of asymmetric reaction. For example, Noyori asymmetric hydrogenation with BINAP/Ru requires a β-ketone, although another catalyst, BINAP/diamine-Ru, widens the scope to α,β-alkenes and aromatic chemicals.
Chiral auxiliaries
A chiral auxiliary is an organic compound which couples to the starting material to form a new compound which can then undergo diastereoselective reactions via intramolecular asymmetric induction. At the end of the reaction the auxiliary is removed, under conditions that will not cause racemization of the product. It is typically then recovered for future use.
Chiral auxiliaries must be used in stoichiometric amounts to be effective and require additional synthetic steps to append and remove the auxiliary. However, in some cases the only available stereoselective methodology relies on chiral auxiliaries and these reactions tend to be versatile and very well-studied, allowing the most time-efficient access to enantiomerically pure products. Additionally, the products of auxiliary-directed reactions are diastereomers, which enables their facile separation by methods such as column chromatography or crystallization.
Biocatalysis
Biocatalysis makes use of biological compounds, ranging from isolated enzymes to living cells, to perform chemical transformations.
The advantages of these reagents include very high e.e.s and reagent specificity, as well as mild operating conditions and low environmental impact. Biocatalysts are more commonly used in industry than in academic research; for example in the production of statins.
The high reagent specificity can be a problem, however, as it often requires that a wide range of biocatalysts be screened before an effective reagent is found.
Enantioselective organocatalysis
Organocatalysis refers to a form of catalysis, where the rate of a chemical reaction is increased by an organic compound consisting of carbon, hydrogen, sulfur and other non-metal elements.
When the organocatalyst is chiral, then enantioselective synthesis can be achieved;
for example a number of carbon–carbon bond forming reactions become enantioselective in the presence of proline with the aldol reaction being a prime example.
Organocatalysis often employs natural compounds and secondary amines as chiral catalysts; these are inexpensive and environmentally friendly, as no metals are involved.
Chiral pool synthesis
Chiral pool synthesis is one of the simplest and oldest approaches for enantioselective synthesis. A readily available chiral starting material is manipulated through successive reactions, often using achiral reagents, to obtain the desired target molecule. This can meet the criteria for enantioselective synthesis when a new chiral species is created, such as in an SN2 reaction.
Chiral pool synthesis is especially attractive for target molecules having similar chirality to a relatively inexpensive naturally occurring building-block such as a sugar or amino acid. However, the number of possible reactions the molecule can undergo is restricted and tortuous synthetic routes may be required (e.g. Oseltamivir total synthesis). This approach also requires a stoichiometric amount of the enantiopure starting material, which can be expensive if it is not naturally occurring.
Separation and analysis of enantiomers
The two enantiomers of a molecule possess many of the same physical properties (e.g. melting point, boiling point, polarity etc.) and so behave identically to each other. As a result, they will migrate with an identical Rf in thin layer chromatography and have identical retention times in HPLC and GC. Their NMR and IR spectra are identical.
This can make it very difficult to determine whether a process has produced a single enantiomer (and crucially which enantiomer it is) as well as making it hard to separate enantiomers from a reaction which has not been 100% enantioselective. Fortunately, enantiomers behave differently in the presence of other chiral materials and this can be exploited to allow their separation and analysis.
Enantiomers do not migrate identically on chiral chromatographic media, such as quartz or standard media that has been chirally modified. This forms the basis of chiral column chromatography, which can be used on a small scale to allow analysis via GC and HPLC, or on a large scale to separate chirally impure materials. However this process can require large amount of chiral packing material which can be expensive. A common alternative is to use a chiral derivatizing agent to convert the enantiomers into a diastereomers, in much the same way as chiral auxiliaries. These have different physical properties and hence can be separated and analysed using conventional methods. Special chiral derivitizing agents known as 'chiral resolution agents' are used in the NMR spectroscopy of stereoisomers, these typically involve coordination to chiral europium complexes such as Eu(fod)3 and Eu(hfc)3.
The separation and analysis of component enantiomers of a racemic drugs or pharmaceutical substances are referred to as chiral analysis. or enantioselective analysis. The most frequently employed technique to carry out chiral analysis involves separation science procedures, specifically chiral chromatographic methods.
The enantiomeric excess of a substance can also be determined using certain optical methods. The oldest method for doing this is to use a polarimeter to compare the level of optical rotation in the product against a 'standard' of known composition. It is also possible to perform ultraviolet-visible spectroscopy of stereoisomers by exploiting the Cotton effect.
One of the most accurate ways of determining the chirality of compound is to determine its absolute configuration by X-ray crystallography. However this is a labour-intensive process which requires that a suitable single crystal be grown.
History
Inception (1815–1905)
In 1815 the French physicist Jean-Baptiste Biot showed that certain chemicals could rotate the plane of a beam of polarised light, a property called optical activity.
The nature of this property remained a mystery until 1848, when Louis Pasteur proposed that it had a molecular basis originating from some form of dissymmetry,
with the term chirality being coined by Lord Kelvin a year later.
The origin of chirality itself was finally described in 1874, when Jacobus Henricus van 't Hoff and Joseph Le Bel independently proposed the tetrahedral geometry of carbon. Structural models prior to this work had been two-dimensional, and van 't Hoff and Le Bel theorized that the arrangement of groups around this tetrahedron could dictate the optical activity of the resulting compound through what became known as the Le Bel–van 't Hoff rule.
In 1894 Hermann Emil Fischer outlined the concept of asymmetric induction; in which he correctly ascribed selective the formation of D-glucose by plants to be due to the influence of optically active substances within chlorophyll. Fischer also successfully performed what would now be regarded as the first example of enantioselective synthesis, by enantioselectively elongating sugars via a process which would eventually become the Kiliani–Fischer synthesis.
The first enantioselective chemical synthesis is most often attributed to Willy Marckwald, Universität zu Berlin, for a brucine-catalyzed enantioselective decarboxylation of 2-ethyl-2-methylmalonic acid reported in 1904. A slight excess of the levorotary form of the product of the reaction, 2-methylbutyric acid, was produced; as this product is also a natural product—e.g., as a side chain of lovastatin formed by its diketide synthase (LovF) during its biosynthesis—this result constitutes the first recorded total synthesis with enantioselectivity, as well other firsts (as Koskinen notes, first "example of asymmetric catalysis, enantiotopic selection, and organocatalysis"). This observation is also of historical significance, as at the time enantioselective synthesis could only be understood in terms of vitalism. At the time many prominent chemists such as Jöns Jacob Berzelius argued that natural and artificial compounds were fundamentally different and that chirality was simply a manifestation of the 'vital force' which could only exist in natural compounds. Unlike Fischer, Marckwald had performed an enantioselective reaction upon an achiral, un-natural starting material, albeit with a chiral organocatalyst (as we now understand this chemistry).
Early work (1905–1965)
The development of enantioselective synthesis was initially slow, largely due to the limited range of techniques available for their separation and analysis.
Diastereomers possess different physical properties, allowing separation by conventional means, however at the time enantiomers could only be separated by spontaneous resolution (where enantiomers separate upon crystallisation) or kinetic resolution (where one enantiomer is selectively destroyed). The only tool for analysing enantiomers was optical activity using a polarimeter, a method which provides no structural data.
It was not until the 1950s that major progress really began. Driven in part by chemists such as R. B. Woodward and Vladimir Prelog but also by the development of new techniques.
The first of these was X-ray crystallography, which was used to determine the absolute configuration of an organic compound by Johannes Bijvoet in 1951.
Chiral chromatography was introduced a year later by Dalgliesh, who used paper chromatography to separate chiral amino acids.
Although Dalgliesh was not the first to observe such separations, he correctly attributed the separation of enantiomers to differential retention by the chiral cellulose. This was expanded upon in 1960, when Klem and Reed first reported the use of chirally-modified silica gel for chiral HPLC separation.
Thalidomide
While it was known that the different enantiomers of a drug could have different activities, with significant early work being done by Arthur Robertson Cushny, this was not accounted for in early drug design and testing. However, following the thalidomide disaster the development and licensing of drugs changed dramatically.
First synthesized in 1953, thalidomide was widely prescribed for morning sickness from 1957 to 1962, but was soon found to be seriously teratogenic, eventually causing birth defects in more than 10,000 babies. The disaster prompted many countries to introduce tougher rules for the testing and licensing of drugs, such as the Kefauver-Harris Amendment (US) and Directive 65/65/EEC1 (EU).
Early research into the teratogenic mechanism, using mice, suggested that one enantiomer of thalidomide was teratogenic while the other possessed all the therapeutic activity. This theory was later shown to be incorrect and has now been superseded by a body of research. However it raised the importance of chirality in drug design, leading to increased research into enantioselective synthesis.
Modern age (since 1965)
The Cahn–Ingold–Prelog priority rules (often abbreviated as the CIP system) were first published in 1966; allowing enantiomers to be more easily and accurately described.
The same year saw first successful enantiomeric separation by gas chromatography an important development as the technology was in common use at the time.
Metal-catalysed enantioselective synthesis was pioneered by William S. Knowles, Ryōji Noyori and K. Barry Sharpless; for which they would receive the 2001 Nobel Prize in Chemistry. Knowles and Noyori began with the development of asymmetric hydrogenation, which they developed independently in 1968. Knowles replaced the achiral triphenylphosphine ligands in Wilkinson's catalyst with chiral phosphine ligands. This experimental catalyst was employed in an asymmetric hydrogenation with a modest 15% enantiomeric excess. Knowles was also the first to apply enantioselective metal catalysis to industrial-scale synthesis; while working for the Monsanto Company he developed an enantioselective hydrogenation step for the production of L-DOPA, utilising the DIPAMP ligand.
Noyori devised a copper complex using a chiral Schiff base ligand, which he used for the metal–carbenoid cyclopropanation of styrene.
In common with Knowles' findings, Noyori's results for the enantiomeric excess for this first-generation ligand were disappointingly low: 6%. However continued research eventually led to the development of the Noyori asymmetric hydrogenation reaction.
Sharpless complemented these reduction reactions by developing a range of asymmetric oxidations (Sharpless epoxidation, Sharpless asymmetric dihydroxylation, Sharpless oxyamination) during the 1970s and 1980s. With the asymmetric oxyamination reaction, using osmium tetroxide, being the earliest.
During the same period, methods were developed to allow the analysis of chiral compounds by NMR; either using chiral derivatizing agents, such as Mosher's acid,
or europium based shift reagents, of which Eu(DPM)3 was the earliest.
Chiral auxiliaries were introduced by E.J. Corey in 1978 and featured prominently in the work of Dieter Enders. Around the same time enantioselective organocatalysis was developed, with pioneering work including the Hajos–Parrish–Eder–Sauer–Wiechert reaction.
Enzyme-catalyzed enantioselective reactions became more and more common during the 1980s, particularly in industry, with their applications including asymmetric ester hydrolysis with pig-liver esterase. The emerging technology of genetic engineering has allowed the tailoring of enzymes to specific processes, permitting an increased range of selective transformations. For example, in the asymmetric hydrogenation of statin precursors.
| Physical sciences | Synthetic strategies | Chemistry |
1464343 | https://en.wikipedia.org/wiki/Agricultural%20extension | Agricultural extension | Agricultural extension is the application of scientific research and new knowledge to agricultural practices through farmer education. The field of 'extension' now encompasses a wider range of communication and learning activities organized for rural people by educators from different disciplines, including agriculture, agricultural marketing, health, and business studies.
Extension practitioners can be found throughout the world, usually working for government agencies. They are represented by several professional organizations, networks and extension journals.
Agricultural extension agencies in developing countries receive large amounts of support from international development organizations such as the World Bank and the Food and Agriculture Organization of the United Nations.
Extension terminology
The use of the word 'extension' originated in England in 1866. Modern extension began in Dublin, Ireland in 1847 with Lord Clarendon's itinerant instructors during the Great Famine. It expanded in Germany in the 1850s, through the itinerant agricultural teachers called Wanderlehrer, and later in the United States via the cooperative extension system authorized by the Smith-Lever Act of 1914. The term was later adopted in the USA, while in Britain it was replaced with "advisory service" in the 20th century. A number of other terms are used in different parts of the world to describe the same or similar concept:
Arabic: Al-Ershad ("guidance")
Bengali: সম্প্রসারণ (shomprosharon 'extension')
Dutch: Voorlichting ("lighting the path")
German: Beratung ("advisory work")
French: Vulgarisation ("popularization")
Italian: Assistenza tecnica e divulgazione agricola ("agricultural technical assistance and popularization")
Spanish: Capacitación ("training", "capacity building")
Thai, Lao: Song-Suem ("to promote")
Persian: Tarvij & Gostaresh ("to promote and to extend") - ترویج و گسترش
Hindi: vistaar ("to extend")- विस्तार
Somali: hormarin & ballaarin' ("to promote and extend"')
In the US, an extension agent is a university employee who develops and delivers educational programs to assist people in economic and community development, leadership, family issues, agriculture and environment. Another program area provided by extension agents is 4-H and youth activities. Many extension agents work for cooperative extension service programs at land-grant universities. They are sometimes referred to as county agents, or extension educators. Often confused with extension agents, extension specialists are subject-matter experts usually employed as scientists and university professors in various departments in the land-grant university system. Subjects range from agriculture, life sciences, economics, engineering, food safety, pest management, veterinary medicine, and various other allied disciplines. These subject matter specialists work with agents (usually in a statewide or regional team environment) to support programs within the cooperative extension system.
Historical definitions
The examples given below are taken from a number of books on extension published over a period of more than 50 years:
1949: The central task of extension is to help rural families help themselves by applying science, whether physical or social, to the daily routines of farming, homemaking, and family and community living.
1965: Agricultural extension has been described as a system of out-of-school education for rural people.
1966: Extension personnel have the task of bringing scientific knowledge to farm families in the farms and homes. The object of the task is to improve the efficiency of agriculture.
1973: Extension is a service or system which assists farm people, through educational procedures, in improving farming methods and techniques, increasing production efficiency and income, bettering their standard of living and lifting social and educational standards.
1974: Extension involves the conscious use of communication of information to help people form sound opinions and make good decisions.
1982: Agricultural Extension: Assistance to farmers to help them identify and analyze their production problems and become aware of the opportunities for improvement.
1988: Extension is a professional communication intervention deployed by an institution to induce change in voluntary behaviors with a presumed public or collective utility.
1997: Extension is the organized exchange of information and the deliberate transfer of skills.
1999: The essence of agricultural extension is to facilitate interplay and nurture synergies within a total information system involving agricultural research, agricultural education and a vast complex of information-providing businesses.
2004: Extension is a series of embedded communicative interventions that are meant, among other goals, to develop and/or induce innovations which help to resolve (usually multi-actor) problematic situations.
2006: Extension is the process of enabling change in individuals, communities and industries involved in the primary industry sector and in natural resource management.
History
Origins of agricultural extension
It is not known where or when the first extension activities took place. It is known, however, that Chinese officials were creating agricultural policies, documenting practical knowledge, and disseminating advice to farmers at least 2,000 years ago. For example, in approximately 800 BC, the minister responsible for agriculture under one of the Zhou dynasty emperors organized the teaching of crop rotation and drainage to farmers. The minister also leased equipment to farmers, built grain stores and supplied free food during times of famine.
The birth of the modern extension service has been attributed to events that took place in Ireland in the middle of the 19th century. Between 1845–51 the Irish potato crop was destroyed by fungal diseases and a severe famine occurred. The British Government arranged for "practical instructors" to travel to rural areas and teach small farmers how to cultivate alternative crops. This scheme attracted the attention of government officials in Germany, who organized their own system of traveling instructors. By the end of the 19th century, the idea had spread to Denmark, Netherlands, Italy, and France.
The term "university extension" was first used by the Universities of Cambridge and Oxford in 1867 to describe teaching activities that extended the work of the institution beyond the campus. Most of these early activities were not, however, related to agriculture. It was not until the beginning of the 20th century, when colleges in the United States started conducting demonstrations at agricultural shows and giving lectures to farmer’s clubs, that the term "extension service" was applied to the type of work that we now recognize by that name.
In the United States, the Hatch Act of 1887 established a system of agricultural experiment stations in conjunction with each state's land-grant university, and the Smith-Lever Act of 1914 created a system of cooperative extension to be operated by those universities in order to inform people about current developments in agriculture, home economics, and related subjects.
In 1966 the National Sea Grant College Program was established and Bill Wick developed the first Marine Advisory Program in Oregon using the Extension model. The first Marine Extension agent was Bob Jacobsen, and was known as "an agricultural agent in hip-boots".
Four generations of extension in Asia
The development of extension services in modern Asia has differed from country to country. Despite the variations, it is possible to identify a general sequence of four periods or "generations":
Colonial agriculture: Experimental stations were established in many Asian countries by the colonial powers. The focus of attention was usually on export crops such as rubber, tea, cotton, and sugar. Technical advice was provided to plantation managers and large landowners. Assistance to small farmers who grew subsistence crops was rare, except in times of crisis.
Diverse top-down extension: After independence, commodity-based extension services emerged from the remnants of the colonial system, with production targets established as part of five-year development plans. In addition, various schemes were initiated to meet the needs of small farmers, with support from foreign donors.
Unified top-down extension: During the 1970s and 1980s, the Training and Visit system (T&V) was introduced by the World Bank. Existing organizations were merged into a single national service. Regular messages were delivered to groups of farmers, promoting the adoption of "Green Revolution" technologies.
Diverse bottom-up extension: When World Bank funding came to an end, the T&V system collapsed in many countries, leaving behind a patchwork of programs and projects funded from various other sources. The decline of central planning, combined with a growing concern for sustainability and equity, has resulted in participatory methods gradually replacing top-down approaches.
The fourth generation is well established in some countries, while it has only just begun in other places. While it seems likely that participatory approaches will continue to spread in the next few years, it is impossible to predict the long-term future of extension. Compared to 20 years ago, agricultural extension now receives considerably less support from donor agencies. Among academics working in this field, some have recently argued that agricultural extension needs to be reinvented as a professional practice. Other authors have abandoned the idea of extension as a distinct concept and prefer to think in terms of "knowledge systems" in which farmers are seen as experts rather than adopters.
Aspects of future extension education:
Evolution of extension system and operationalisation of approaches
Future extension education initiatives
Collegiate participation of farmers
Web enabled technology dissemination
Developing cases as tools for technology dissemination
Agriculture as a profitable venture
Scaling up of group mobilization
Micro-enterprises promotion
Several of the institutional innovations that have come up in response to the weaknesses in public research and extension system have given enough indications of the emergence of an agricultural innovation system in India. This has resulted in the blurring of the clearly demarcated institutional boundaries between research, extension, farmers, farmers' groups, NGOs and private enterprises. Extension should play the role of facilitating the access to and transfer of knowledge among the different entities involved in the innovation system and create competent institutional modes to improve the overall performance of the innovation system. Inability to play this important role would further marginalize extension efforts.
Communication processes
The term "extension" has been used to cover widely differing communication systems. Two particular issues help to define the type of extension: how communication takes place, and why it takes place.
The related but separate field of agricultural communication has emerged to contribute to in-depth examinations of the communication processes among various actors within and external to the agricultural system. This field refers to the participatory extension model as a form of public relations-rooted two-way symmetric communication based on mutual respect, understanding, and influence between an organization and its stakeholders.
Agricultural communication can take three modes—face-to-face training, training "products" such as manuals and videos, or information and communication technologies (ICTs), such as radio and short message system (SMS). The most effective systems facilitate two-way communication and often combine different modes.
Four paradigms of agricultural extension
Any particular extension system can be described in terms of both how communication takes place and why it takes place. It is not the case that paternalistic systems are always persuasive, nor is it the case that participatory projects are necessarily educational. Instead, there are four possible combinations, each of which represents a different extension paradigm, as follows:
Technology transfer (persuasive + paternalistic): This paradigm was prevalent in colonial times and reappeared in the 1970s and 1980s when the "Training and Visit" system was established across Asia. Technology transfer involves a top-down approach that delivers specific recommendations to farmers about the practices they should adopt.
Advisory work (persuasive + participatory): This paradigm can be seen today where government organizations or private consulting companies respond to farmers' inquiries with technical prescriptions. It also takes the form of projects managed by donor agencies and NGOs that use participatory approaches to promote predetermined packages of technology.
Human resource development (educational + paternalistic): This paradigm dominated the earliest days of extension in Europe and North America, when universities gave training to rural people who were too poor to attend full-time courses. It continues today in the outreach activities of colleges around the world. Top-down teaching methods are employed, but students are expected to make their own decisions about how to use the knowledge they acquire.
Facilitation for empowerment (educational + participatory): This paradigm involves methods such as experiential learning and farmer-to-farmer exchanges. Knowledge is gained through interactive processes and the participants are encouraged to make their own decisions. The best known examples in Asia are projects that use Farmer Field Schools (FFS) or participatory technology development (PTD).
There is some disagreement about whether or not the concept and name of 'extension' really encompasses all four paradigms. Some experts believe that the term should be restricted to persuasive approaches, while others believe it should only be used for educational activities. Paulo Freire has argued that the terms ‘extension’ and ‘participation’ are contradictory. There are philosophical reasons behind these disagreements. From a practical point of view, however, communication processes that conform to each of these four paradigms are currently being organized under the name of extension in one part of the world or another. Pragmatically, if not ideologically, all of these activities are considered to be represented in agricultural extension.
| Technology | Academic disciplines | null |
1465291 | https://en.wikipedia.org/wiki/Squat%20lobster | Squat lobster | Squat lobsters are dorsoventrally flattened crustaceans with long tails held curled beneath the cephalothorax. They are found in the two superfamilies Galatheoidea and Chirostyloidea, which form part of the decapod infraorder Anomura, alongside groups including the hermit crabs and mole crabs. They are distributed worldwide in the oceans, and occur from near the surface to deep sea hydrothermal vents, with one species occupying caves above sea level. More than 900 species have been described, in around 60 genera. Some species form dense aggregations, either on the sea floor or in the water column, and a small number are commercially fished.
Description
The two main groups of squat lobsters share most features of their morphology. They resemble true lobsters in some ways, but are somewhat flattened dorsoventrally, and are typically smaller, ranging from 0.7 to 3.5 inches in length. Squat lobsters vary in postorbital carapace length (measured from the eye socket to the rear edge), from in the case of Munidopsis aries, down to only a few millimetres in the case of Galathea intermedia and some species of Uroptychus. As in other decapod crustaceans, the bilaterally symmetrical body of a squat lobster may be divided into two main regions: the cephalothorax (itself made up of the cephalon, or head, and the thorax), and the pleon or abdomen. The pleon only being partly flexed under the cephalothorax and the cephalothorax being more long than it is broad makes the squat lobster a morphological intermediate between a lobster and crab.
The cephalothorax is made of 19 body segments (somites), although the divisions are not obvious and are most easily inferred from the paired appendages. From front to back, these are the two pairs of antennae, six pairs of mouthparts (mandibles, maxillae, maxillules and three pairs of maxillipeds), and five pairs of pereiopods. The cephalothorax is covered with a thick carapace, which may extend forwards in front of the eyes to form a rostrum; this is highly variable among squat lobsters, being vestigial in Chirostylus, wide and often serrated in some genera, and long, narrow, and flanked with "supraorbital spines" in others. The degree of ornamentation on the surface of the carapace also varies widely, and there are almost always at least a few setae (bristles), which can be iridescent in some members of the Galatheidae and Munididae. A pair of compound eyes also project on stalks from the front of the carapace; these are made up of ommatidia with square facets, which is typical of the "reflecting superposition" form of eye. Many deep-sea species have reduced eyes, and reduced movement of the eyestalks. In the families Munididae and Galatheidae, there is often a row of setae close to the eyes, forming "eyelashes".
The mouthparts consist of six pairs of appendages— three posterior cephalic appendages and the first three pairs of thoracic appendages. While their function was traditionally believed to be limited to food handling, the mouthparts have a more complex movement pattern that allows them to perform a variety of functions such as prey- and sediment-gathering, sediment transfer, and sediment sorting/particle rejection. In Munida Sarsi, the farther the mouthparts are located from the mouth, the more complex in movement and functional scheme they are.
The most conspicuous appendages are the pereiopods, and the largest of these is the first pair. These each end in a chela (claw), and are therefore known as the "chelipeds"; they can be more than six times the body length, although some groups show sexual dimorphism, with females having proportionally shorter chelipeds. The following three pairs of pereiopods are somewhat smaller than the chelipeds and are without claws, but are otherwise similar; they are used for walking. The fifth pair of pereiopods are much smaller than the preceding pairs, and are held inconspicuously under the carapace. They each end in a tiny chela, and are generally believed to be used for cleaning the body, especially the gills, which are in a cavity protected by the carapace.
The pleon is made up of six somites, each bearing a pair of pleopods, and terminating in a telson. The first somite is narrower than the succeeding somites, and the last pair of pleopods are modified into uropods, which flank the telson. The pleon is usually curled under the thorax, such that only the first three somites are visible from above. The form of the pleopods varies between the sexes. In females, the first one or two pairs are missing, while the remaining pairs are uniramous, and have long setae, to which the eggs can be attached. In males, the first two pairs are formed into gonopods, and are used to transfer the spermatophore to the female during mating; the first pair is often missing. The remaining pleopods can be similar to those of the females, or reduced in size, or entirely absent. In both sexes, the uropods are biramous.
Carcinisation has previously been explored in regards to outer morphology; however, the external change in body shaped has influence on the internal anatomical features as well. The use of micro-computer tomography and 3D reconstruction have brought to light anatomical disparity within Galatheoidea. Differences have been found in the ventral vessel system between porcelain crabs and squat lobsters. Carcinisation is also responsible for the loss of the caridoid escape reaction which caused a shift in gonads and the pleonal neuromeres for squat lobsters.
Development
Fecundity or number of eggs increases with smaller sized eggs and increasing body size of the parent. This results in increasing incubation time and consequently, increased egg volume. The trend of larger numbers of eggs and smaller sized eggs is mostly found in lower latitudes and cooler temperatures in order to accommodate for the longer incubation time.
The development period of the embryo consists of five distinct stages in which the lasts several months. Throughout the five stages, both the diameter and volume of the egg increases. In the first stage of embryo development, the spherical egg is a uniform dark color. In stage II, the optic lobe and appendages begin formation. Stage III is when the abdominal segments and terminal spines begin to develop. In stage IV, there is pigmentation of the ocular lobe, segmentation of the maxillipeds, and cardiac movement. In the fifth and final stage, the eyes are enlarged and the abdomen is extended.
Ecology and behavior
Aggregation and migration
Benthic aggregation
Squat lobster aggregation is theorized to be proportional to the amount of available organic particulate carbon reaching the seafloor. As such, many species of benthic squat lobsters aggregate into groups of very high population density around a number of different types of highly productive areas of the deep sea, like hydrothermal vents, cold seeps, sites of food falls like whale falls or wood falls, and shipwrecks. Squat lobsters seem to aggregate in these areas because they are associated with the fauna that thrives in them, like Bathymodiolus mussels and vestimentiferan tubeworms, but they are also attracted to the three-dimensional structures found in them. Burrows, crevices, or debris from shipwrecks are suggested to serves as shelter against predators, as well as food accumulation sites. A variety of squat lobster called Emmunida picta were found to almost exclusively reside on some type of structure, mostly Lophelia pertusa.
Squat lobster species found on seamounts typically have smaller bodies with shorter larval stages, as opposed to rise and ridge habitats. It has been suggested that this is due to the difference in substrates at these habitats.
Ontogenetic migrations and pelagic aggregation
Pleuroncodes planipes perform vertical migration into the water column from the benthos at different stages in development, an example of ontogenetic niche shift. Early larval stages are found mostly near the sea surface, but older larval and juvenile stages have a wide vertical distribution through the water column, and adults become almost purely benthic, with only a few vertical migrations in their remaining lifetime. In these migrations, they will form large swarms (up to 200 m vertically, and up to 10 km horizontally) in which they rely the selective tidal stream for feeding.
Munida gregaria form aggregations in warm summer waters of the Pacific Ocean associated with river plume fronts, headland fronts, and shallow internal waves. Density of these aggregations are, on average, 2700 individuals per cubic meter. M. gregaria are able to aggregate in the pelagic region due to a number of unique features as compared to benthic squat lobsters, including fast swimming speeds, reduced density, reduced sinking rates as a result of greater morphological surface area, and optimized aerobic metabolism. M. gregaria also exhibit ontogenetic migration through larvae accumulation in highly productive nearshore waters, which then move toward the mid-continental shelf as they mature, and move completely offshore around full maturation.
In 2020, a study of squat lobsters determined that these crustaceans are far more diverse than previously thought. Through this study, 16 new species within the Leiogalathea genera were described. It was also revealed that diversity of squat lobsters in the Atlantic Ocean is relatively poor in comparison with the Pacific Ocean.
Claw position behavior
In a resting posture, squat lobsters rest their claws on the substrate in front of them. E. picta were observed most frequently in all conditions with their claws extended into the water column, perpendicular to the substrate. In the case of this study, the behavior was thought to be an avoidance response to the surveillance submersible, hence why this behavior was so often observed. In general, it is thought that this behavior may be a mechanism to increase the perceived size of the squat lobster as an aggressive or perhaps illusory display to ward off predators, as well as an active "fishing" strategy to catch prey.
Aggression and agonistic behavior
While squat lobsters look like true lobsters, they are more closely related to hermit crabs. Instead of carrying shells on their backs, they squeeze their bodies into crevices and leave their claws exposed to defend themselves from predators or other squat lobsters.
Squat lobsters are generally unaggressive toward each other, but instances can occur in particular scenarios. Individuals among dense populations will make decisions about whether to hunt for food or engage in deposit feeding on the basis of minimizing aggressive interactions. In general, squat lobsters exhibit no lasting dominance hierarchies, nor do they engage in territorialist behavior. When aggressive displays do occur, as a result of competition for mates or food, the aggressive behavior is ignored 70% of the time, met with submission 20% of the time, and met with reciprocal aggression 10% of the time. A 2001 study examined the effects of serotonergic and octopaminergic systems in Munida quadrispina, and found that injected serotonin elicits aggressive postures and behaviors, including increased likelihood and intensity of aggressive reactions to real or artificial squat lobsters, while injected octopamine reduced instances of aggressive behavior, including increasing the likelihood of escape responses.
Feeding
Squat lobsters feed on a variety of foods, with some species filter feeding, while others are detritus-feeders, algal grazers, scavengers, predators, and occasionally cannibals. Some are highly specialised; Munidopsis andamanica is a deep-sea species that feeds only on sunken wood, including trees washed out to sea and timber from shipwrecks. Squat lobsters are large enough to be caught by top predators, and can thus form a "direct trophic shortcut" between the primary producers at the bottom of the food web, and the carnivores at the top.
Breeding
Squat lobsters, in particular Cervimunida johni and Pleuroncodes monodon, are known to mate during the intermolt period. Squat lobster mating was shown to be unrelated to the female molt period, unlike many other crustacean species; instead, mating occurs when most females are found to be ovigerous and there is no or only minor molting activity, i.e the intermolt period. These species of squat lobster also displayed very short interbrood intervals, or time between mating, generally not longer than a few days.
Fisheries
Flesh from these animals is often commercially sold in restaurants as "langostino" or sometimes dishonestly called "lobster" when incorporated in seafood dishes. As well as being used for human consumption, there is demand for squat lobster meat to be used as feed in fish farms and shrimp or prawn farms. This is in part because they contain astaxanthin, a pigment that helps to colour the meat of farmed salmon and trout.
Despite their worldwide distribution and great abundance, there are few functioning fisheries for squat lobsters. Experimental fisheries have occurred in several countries, including Argentina, Mexico, and New Zealand, but commercial exploitation is restricted to Latin America, and chiefly to Chile. The main target species are Pleuroncodes monodon, P. planipes, and Cervimunida johni.
In Central America, the primary species of squat lobster targeted by fisheries is a species of Pleuroncodes. There is a great deal of confusion over both scientific names and common names, and the exact species is often unknown. In El Salvador, for instance, the commercial catch is generally referred to as "P. planipes", but is in fact P. monodon. Commercial fishing for squat lobsters in El Salvador began in the early 1980s; production increased markedly in the 2001 season, and has continued to grow, now [when?] making up 98% of the demersal resources landed in El Salvador, with annual catches peaking at 13,708 t in 2005. In Costa Rica, aggregations of squat lobsters are avoided, as the fishermen fear the squat lobsters will clog their nets. In Nicaragua, squat lobsters are heavily exploited, especially following a large increase in fishing effort in the 2007 season. In Panama, production reached 492 t in 2008. Chilean squat lobster fisheries initially targeted Cervimunida johni, beginning in 1953. By the mid-1960s, effort had largely switched to P. monodon. In an effort to conserve stocks, the Chilean government instituted quotas for squat lobsters, and the fishery is closely monitored. In New Zealand, Munida gregaria has been considered as a potential fisheries resource, particularly to feed farmed Chinook salmon (Oncorhynchus tshawytscha).
Classification
Broadly, squat lobsters are classified into two superfamilies: Chirostyloidea and Galatheoidea. Chirostyloidea contain the families Chirostylidae, Eumunididae, and Kiwaidae. Galatheoidea contain the families Galatheidae, Munididae, Munidopsidae, and Porcellanidae. The systematics of squat lobsters and the classification of deep-sea squat lobsters is an area of active research due to the limited fossil record. Deep-sea squat lobsters, existing at depths greater than 200m, are classified in the Mundidiae and Munidopsidae families of the superfamily Galatheoidea. Deep-sea squat lobsters display greater morphological divergences and lower genetic divergences evolutionary compared to their shallow-water counterparts. In early classifications, squat lobsters were placed in the superfamily Galatheoidea alongside the porcelain crabs of the family Porcellanidae. This relationship, however, was re-examined in the late 21st-century. Molecular and morphological data indicate that Galatheoidea is not a monophyletic group; Galatheidae, Porcellanidae, Kiwaidae, and Chirostylidae have independent origins. Galatheoidea is the largest superfamily in the Anomura suborder, which is reflected in the wide range of habitats of species belonging to this "superfamily." Few morphological characteristics distinguish squat lobsters from other families in the Anomura. Hence, deep-sea squat lobsters were classified in the Chirostyloidea superfamily in 2012 as DNA sequencing indicated that squat lobsters are not a monophyletic group. For example, Chirostylidae and Kiwaidae are distantly related to the other squat lobsters, and are closer related to hermit crabs and king crabs (Paguroidea), the mole crabs in the superfamily Hippoidea, and the small families Lomisidae and Aeglidae. Squat lobsters continue to be described in both the Chirostyloidea and Galatheoidea superfamilies.
Evolutionary history
Squat lobsters contain a total of around 60 genera, divided into over 900 recognized species; more than 120 undescribed species likely exist. It is likely that squat lobsters underwent deep sea colonization multiple times in evolutionary history. Squat lobsters underwent rapid diversification in the late Oligocene through the Miocene likely due to a variety of selection pressures, beginning in the Southwest Pacific. Fossil galatheoid squat lobsters have been found in strata dating back to the Middle Jurassic of Europe. No fossils are currently assigned to the Chirostyloidea. Pristinaspina may belong either in the family Kiwaidae or Chirostylidae.
Distribution
Generally, species richness of deep-sea squat lobsters increases with proximity to the equator and the Western Pacific. A split in the galatheid fauna exists between the Eastern Pacific, containing much less species richness, and the Western Pacific. The center of diversity for squat lobsters is the "coral triangle", or Indo-Australian Archipelago, especially in the region of New Caledonia (with more than 300 species) and the region of Indonesia and the Philippines. High endemism is reported in this area and it is hypothesized this is due to many independent evolutions. This region results in high diversification due to global warming, tectonic activity, and oceanic currents. Modern regions of high speciation of squat lobsters includes seamounts near ocean trenches in the West Pacific, but squat lobsters are found in a range of habitats including continental shelfs, ridges, and abyssal seabeds. The hotspot distribution of squat lobsters in shallow waters seems to mirror the distribution of deep-sea species.
In March 2022 it was reported that a squat lobster, possibly from the genus Munidopsis, had been filmed on the wreck of the Endurance, which sank in 1915 in the Antarctic. This was the first record of a living squat lobster in the Weddell Sea.
| Biology and health sciences | Crabs and hermit crabs | Animals |
1465643 | https://en.wikipedia.org/wiki/Oslo%20Metro | Oslo Metro | The Oslo Metro ( or or simply ) is the rapid transit system of Oslo, Norway, operated by Sporveien T-banen on contract from the transit authority Ruter. The network consists of five lines that all run through the city centre, with a total length of , serving 101 stations of which 17 are underground or indoors. In addition to serving 14 out of the 15 boroughs of Oslo, two lines run to Kolsås and Østerås, in the neighbouring municipality of Bærum. In 2016, the system had an annual ridership of 118 million.
The first rapid transit line, the Holmenkollen Line, opened in 1898, with the branch Røa Line opening in 1912. It became the first Nordic underground rapid transit system in 1928, when the underground line to Nationaltheatret was opened. After 1993 trains ran under the city between the eastern and western networks in the Common Tunnel, followed by the 2006 opening of the Ring Line. All the trains are operated with MX3000 stock. These replaced the older T1000 stock between 2006 and 2010.
History
Suburban lines in the west
Rail transport in Oslo started in 1854, with the opening of Hoved Line to Eidsvoll, through Groruddalen. In 1872, Drammen Line, going through Oslo West, and in 1879, Østfold Line going through Nordstrand opened, offering a limited rail service to those parts of the city. By 1875, Kristiania Sporveisselskab (KSS) opened the first horsecar trams. In 1894 electric trams were in service by Kristiania Elektriske Sporvei (KES).
The first suburban tram line was the Holmenkollen Line that was opened by Holmenkolbanen in 1898; like all the later suburban tram line these were electric trams with a grade-separated right-of-way and proper stations instead of tram stops, making it the first rapid transit in Oslo. Unlike the other suburban tram lines that were built later, the Holmenkollen Line was not extended into the city as a streetcar—instead passengers had to change at Majorstuen to the streetcars, though the system did not take into use wider suburban stock () until 1909. A branch line was opened in 1912, to Smestad, and in 1916 the Holmenkollen Line was extended to Tryvann, with the last part from Frognerseteren single track and used for freight, and removed in 1939.
In 1912, the construction of the first underground railway in the Nordic Countries started, when A/S Holmenkolbanen started construction of an extension of their line from Majorstuen to Nationaltheatret. The line was opened in 1928, with one intermediate station at Valkyrie Plass, giving the two suburban lines access to the central business district of Oslo. In retrospect, it is seen as een early premetro example. It was the second underground railway to be opened in the Nordic countries after Boulevardtunnlen in Copenhagen which opened in July 1918.
The success of the suburban lines tempted KES to extend their streetcar service west from Skøyen as a suburban line; the Lilleaker Line opened to Lilleaker in 1919, to Avløs in 1924 and to Kolsås in 1930. A new section from Jar to Sørbyhaugen opened in 1942, connecting the line from Jar to Kolsås to Nationaltheatret, and making it a rapid transit and the replacement of stock with wide suburban standard. This service remained part of the municipal Oslo Sporveier, that had bought all the streetcar companies in 1924.
Compensation for large amounts of damage to houses along the route during construction, along with higher construction costs than calculated was a heavy burden on the company, and in 1934, the municipality of Aker took over the common stock, though the preferred stock remained listed on the Oslo Stock Exchange until 1975, as Oslo Sporveier gradually took over the operation of the western suburban lines. Akersbanerne opened the connecting Sognsvann Line in 1934.
Metro
The first idea to launch a citywide rapid transit was launched in 1912 with the construction of the Ekeberg Line; constructed with the same width profile as the Holmenkollen Line, the plan was to build a tunnel under the city center and run through trains, but large cost expenditures on the first section of the Common Tunnel ceased the plans. As part of the rebuilding after World War II a planning office for a T-bane was established in 1949, with the first plans launched in 1951; in 1954, the city council decided to build the T-bane network in Eastern Oslo with four branches. The system would feature improvements over the suburban lines in having a third rail power supply, cab signaling with Automatic Train Protection, stations long enough for six-car trains and level crossings replaced by bridges and underpasses—specifications christened metro standard.
At the time there were two suburban tramways on the east side, the Ekeberg Line (opened in 1919) and the Østensjø Line (1923). Only the latter would be connected to the T-bane; the Ekeberg Line would remain a tramway, but three new lines were to be built—the Grorud Line on the north side and the Furuset Line on the south side of Groruddalen and the Lambertseter Line on the east of Nordstrand. These areas were all chosen as new suburbs for Oslo, and would quickly need a good public transport system; suburban lines would first be built out extending from the existing tramway, and later a final section with tunnel to the central station would be built. The Lambertseter Line was opened in 1957, from Brynseng to Bergkrystallen while the Østensjø Line was extended to Bøler in 1958.
The metro opened on 22 May 1966, when the Common Tunnel opened from Brynseng to the new downtown station of Jernbanetorget, located beside the Oslo East Railway Station. In October the Grorud Line opened to Grorud while the Østensjø Line was connected to the system in 1967 when the line also was extended to Skullerud. In 1970, the Furuset Line opened to Haugerud and extended to Trosterud in 1974, at the same time as the Grorud Line was extended to Vestli. By 1981, the Furuset Line had reached Ellingsrudåsen. The metro took delivery of T1000 rolling stock from Strømmens Værksted; from 1964 to 1978, 162 cars in three-car configurations were delivered for the eastern network.
One tunnel
The eastern network was extended from Jernbanetorget to Sentrum in 1977. This station was forced to close in 1983, due to water leakage, and when it opened again in 1987, renamed Stortinget, the west network tunnel had also been extended there. Through services were not possible at the time because of incompatibility of signaling and power equipment. Not until 1993 did the first trains run through the station, after the Sognsvann Line had been rebuilt to "metro standard"; the Røa Line followed in 1995. The Holmenkollen and Kolsås Lines remained non-metro, using dual mode trains that switch to overhead lines at Frøen and Montebello. The western network took delivery of 33 T1300 cars in 1978–81, with an additional 16 converted from T1000. In 1994 twelve T2000 cars were delivered for the Holmenkollen Line.
In 2003 the Ring Line opened, connecting Ullevål stadion to Storo. The following year, construction work caused a tunnel to collapse on the Grorud Line—the system's busiest—forcing a shutdown of the line until December, and creating a havoc of overcrowded replacement buses. In 2006 the ring was completed, to Carl Berners plass. At the same time the Kolsås Line was closed for upgrade to metro standard. In 2003 the section of the Kolsås Line in Bærum closed due to budget disagreements between the two counties; after a year of unpopular replacement buses, the line was reopened, only to close again in 2006 for upgrade to metro standard. Disagreements between the two counties meant the upgrade would be done separately on the two sides of the municipal boundary, with the Oslo side opening first. In 2006 the replacement of existing rolling stock with new MX3000 units commenced. The history of the metro and public transport in Oslo is celebrated at the Oslo Tramway Museum in Majorstuen.
Network
The current route network was introduced on 3 April 2016, with the opening of the connection tunnel from Økern to Sinsen and the new Løren station.
The Oslo Metro operates in all fifteen boroughs of Oslo, as well as reaching a bit inside the neighbouring municipality of Bærum. There are five lines, numbered 1 to 5, each colour-coded. They all pass through the Common Tunnel, serving eight branch lines. In addition two lines operate to the Ring Line. Two branches are served by two lines each: the Grorud branch is served by both lines 4 and 5, while the Lambertseter branch has full-time service by line 4 and limited service by line 1.
The Grorud and Furuset Line head northeast into Groruddalen, while the other two eastern branches head south into Nordstrand. On the west side, the Holmenkoll and Sognsvann Line cover the northern boroughs of Oslo, along with the Ring Line that connects the northeastern and northwestern parts of town. The Kolsås and Røa Line reach deep into the neighbouring municipality of Bærum. All the lines run through the Common Tunnel before reaching out to different lines, or into the Ring. All lines have a base service of four trains per hour while line 2 and the eastern section of line 3 have eight trains per hour weekdays 07:00–19:00. The eastern section of line 2 also has eight trains per hour Saturdays 10:00–19:00. A reduced half-hourly service operates on all lines during early weekend mornings. Trains run from about 05:00 (06:00 at weekends) to 01:00 the next morning.
Lines
Line 1
Line 2
Line 3
Line 4
Line 5
Line 6 (under construction)
A new metro line that will extend from Majorstuen to Fornebu is under construction as of December 2020, aiming to be completed in 2027.
An agreement has been reached and signed between the Oslo city government and the Norwegian state that would share the cost of 13 billion NOK equally between the city and the national government.
Line 7 (proposed)
Stations
The system consists of 101 stations, of which 17 are underground or indoors. The only underground station on the pre-metro western network was Nationaltheatret, and most of the underground station are in the common tunnel under the city center, or in shorter tunnel sections on the eastern network; in particular the Furuset Line runs mainly underground, with all but Haugerud built in or at the opening of a tunnel.
Stations in the city center are located close to large employment centers as well as connection possibilities to other modes of transport, such as tram, rail and bus. All stations can be identified at ground level by signs with a blue T in a circle. Stations outside the center are unmanned since the 1995, with ticket machines for fare purchase; some stations feature kiosks. A system of turnstiles have been installed, but will never be activated due to security issues. All stations have step-free accessibility through at least one entrance (except the inbound platform at Frøen), and the platform height is aligned with the train cars.
Intermodality
The metro is integrated into the public transport system of Oslo and Akershus through the agency Ruter, allowing tickets to also be valid on the Oslo Tramway, city buses, ferries, and the Oslo Commuter Rail operated by Vy. A new, wireless ticketing system, Reisekort, has in the recent years been implemented. As of June 2022, a single ticket for one zone (the entire metro system is in zone 1) costs NOK 39 for adults (a surcharge of 20 NOK is added if you buy onboard within zone 1); a 30-day ticket costs NOK 814 for adults. This includes all means of public transport within the zone where the ticket is first activated (again, for the metro, zone 1). There is a fine of NOK 950, or NOK 1150, for not having a valid ticket, depending on if the fine is paid on location or not.
Oslo maintains a street tram system with six lines, of which two are suburban lines. The street trams operate mostly within the borders of the Ring Line, providing a frequent service in the city centre, with lower average speeds but with more stops. There are major transfer points to the tramway at Majorstuen, Jernbanetorget, Jar, Storo and Forskningsparken.
The commuter train serves suburbs further away from Oslo, though some of the commuter rail services remind of a rapid transit service, in particular line L1 to Lillestrøm and Asker, line L2 to Stabekk and Ski, and line R31 to Jaren with higher service frequency through the continual populated area of Oslo. Transfer to railway services is available at Jernbanetorget (to Oslo S) and Nationaltheatret, the latter with a considerably shorter walk.
Bus services are provided to numerous stations. Most bus services provide feeding to the metro system where possible, and then do not continue into town. However, since the metro operates solely into town, instead of across it, many buses operate between stations on different lines, or provide alternative routes across town.
Future expansion
As part of the political agreement Oslo Package 3, a number of changes and expansions have been proposed for the Oslo Metro. Only one of these, the Fornebu Line, is currently being built - the other projects have been put on hold for it.
Proposed
Expansion of the Furuset Line to Lørenskog with stations at Skårer, Lørenskog Centre and a new terminus at Akershus University Hospital, with travel time to Jernbanetorget of 27 minutes.
A second common tunnel from Majorstuen to Tøyen, creating two new stations at Bislett and southern Grünerløkka (Nybrua). All lines will stop at Stortinget, which will get four platforms. Majorstuen station will be moved underground.
Under construction
The Fornebu Line is planned to run from Majorstuen to the old airport area at Fornebu, through Skøyen and Lysaker; a total of 6 new stations will be on this line. The line began construction in December 2020, with an opening date set around 2029, and will cost an estimated $2.6 billion USD. The line is funded in part by Oslo and Akershus county, among other sources.
Rolling stock
The trains on the Oslo metro are currently exclusively the MX3000, ordered in 2003 to replace the oldest T1000 stock. Delivery started in 2006, and unlike older stock the MX3000 units are painted white instead of red. 83 three-car units were ordered in 2006; a further 32 were ordered in December 2010.
A number of versions of the T1000 stock have earlier been used on the Oslo metro. This includes 146 cars of the types T1 through T4, that have third-rail only operation, and thus did not run on the Holmenkollen and Kolsås lines. These ran usually in units of three or six (sometimes four or five) cars. Types T5 to T8, 49 in total, delivered with both third-rail and overhead wire equipment, normally ran on the Holmenkollen line (two cars) and Kolsås line (three cars).
When the Holmenkollen Line was connected to the T-bane it was still using old teak cars; to allow through services the T2000, capable of dual-system running, was delivered in 1993. They were not particularly successful and only 12 units were delivered, operating in pairs on the Holmenkollen line sometimes connecting with the Lambertseter line, and scrapped in 2010.
Depots and facilities
Avløs Depot – located near Avløs station on the Kolsås Line, it has been closed for refurbishment since 2011 and reopened in August 2015.
Etterstad Depot – located on the shared section of the Østensjø Line, Furuset Line and the Lambertseter Line before Brynseng station, it is used as the main operations centre for the Oslo Metro and has a yard for maintenance of way equipment.
Majorstuen Depot – a small yard used mainly for storing trains, located close to the Oslo Tramway Museum and Majorstuen station.
Ryen Depot – the main storage and maintenance yard for all Oslo Metro trains, located on the Lambertseter Line near Ryen station.
Network map
| Technology | Scandinavia | null |
27900561 | https://en.wikipedia.org/wiki/Livyatan | Livyatan | Livyatan is an extinct genus of macroraptorial sperm whale containing one known species: L. melvillei. The genus name was inspired by the biblical sea monster Leviathan, and the species name by Herman Melville, the author of the famous novel Moby-Dick about a white bull sperm whale. Herman Melville often referred to whales as "Leviathans" in his book. It is mainly known from the Pisco Formation of Peru during the Tortonian stage of the Miocene epoch, about 9.9–8.9 million years ago (mya); however, finds of isolated teeth from other locations such as Chile, Argentina, the United States (California), South Africa and Australia imply that either it or a close relative survived into the Pliocene, around 5mya, and may have had a global presence. It was a member of a group of macroraptorial sperm whales (or "raptorial sperm whales") and was probably an apex predator, preying on whales, seals and so forth. Characteristically of raptorial sperm whales, Livyatan had functional, enamel-coated teeth on the upper and lower jaws, as well as several features suitable for hunting large prey.
Livyatans total length has been estimated to be about , almost similar to that of the modern sperm whale (Physeter macrocephalus), making it one of the largest predators known to have existed. The teeth of Livyatan measured , and are the largest biting teeth of any known animal, excluding tusks. It is distinguished from the other raptorial sperm whales by the basin on the skull spanning the length of the snout. The spermaceti organ contained in that basin is thought to have been used in echolocation and communication, or for ramming prey and other sperm whales. The whale may have interacted with the large extinct shark megalodon (Otodus megalodon), competing with it for a similar food source. Its extinction was probably caused by a cooling event at the end of the Miocene period causing a reduction in food populations. The geological formation where the whale has been found has also preserved a large assemblage of marine life, such as sharks and marine mammals.
Research history
Holotype and naming
In November 2008, a partially preserved skull, as well as teeth and the lower jaw, belonging to L., the holotype specimen MUSM 1676, were discovered in the coastal desert of Peru in the sediments of the Pisco Formation, southwest of the city of Ica. Klaas Post, a researcher for the Natural History Museum Rotterdam in the Netherlands, stumbled across them on the final day of a field trip. The fossils were prepared in Lima, and are now part of the collection of the Museum of Natural History, Lima of National University of San Marcos.
The discoverers originally assigned—in July 2010—the English name of the biblical monster, Leviathan, to the whale as Leviathan melvillei. However, the scientific name Leviathan was also the junior synonym for the mastodon (Mammut), so, in August 2010, the authors rectified this situation by coining a new genus name for the whale, Livyatan, from the original Hebrew name of the monster. The species name melvillei is a reference to Herman Melville, author of the book Moby-Dick, which features a gigantic sperm whale as the main antagonist. The first Livyatan fossils from Peru were initially dated to around 13–12 million years ago (mya) in the Serravallian Age of the Miocene, but this was revised to 9.9–8.9 mya in the Tortonian Age of the Miocene.
Additional specimens
During the late 2010s and 2020s, fossils of large isolated sperm whale teeth were reported from various Miocene and Pliocene localities mostly along the Southern Hemisphere. These teeth have been identified to be of similar size and shape with that of the L. melvillei holotype and may be species of Livyatan. However, it is commonplace that authors do not identify such teeth as a conclusive species of Livyatan, instead opting to assign an open nomenclature in which the biological classifications of the specimens are restricted to comparisons or affinities with Livyatan. This is mostly because isolated teeth tend to not be informative enough to be identified at the species level, meaning that there is some undeterminable possibility that they belong to an undescribed close relative of Livyatan rather than Livyatan itself.
In 2016 in Beaumaris Bay, Australia, a large sperm whale tooth measuring , specimen NMV P16205, was discovered in Pliocene strata by a local named Murray Orr, and was nicknamed the "Beaumaris sperm whale" or the "giant sperm whale". The tooth was donated to Museums Victoria at Melbourne. Though it has not been given a species designation, the tooth looks similar to those of L. melvillei, indicating it was a close relative. The tooth is dated to around 5mya, and so is younger than the L. melvillei holotype by around 4or 5million years.
In 2018, palaeontologists led by David Sebastian Piazza, while revising the collections of the Bariloche Paleontological Museum and the Municipal Paleontological Museum of Lamarque, uncovered two incomplete sperm whale teeth cataloged as MML 882 and BAR-2601 that were recovered from the Saladar Member of the Gran Bajo del Gualicho Formation in the Río Negro Province of Argentina, a deposit that dates between around 20–14 mya. The partial teeth measure and in height, respectively. Anatomical analyses of the specimens found that much of their characteristics are identical to L. melvillei except in width, in which the diameter of both teeth are smaller. Because of this, along with only isolated teeth being available, the palaeontologists chose to assign an open nomenclature, identifying both specimens as aff. Livyatan sp.
In 2019, palaeontologist Romala Govender reported the discovery of two large sperm whale teeth from Pliocene deposits near the Hondeklip Bay village of Namaqualand in South Africa. The pair of teeth, which are stored in the Iziko South African Museum and cataloged as SAM-PQHB-433 and SAM-PQHB-1519, measure and in height, respectively, the latter having its crown missing. Both teeth have open pulp cavities, indicating that both whales were young. The teeth are very similar in shape and size to the mandibular teeth of the L. melvillei holotype, and were identified as cf. Livyatan. Like the Beaumaris specimen, the South African teeth are dated to around 5mya.
In 2023, graduate student Kristin Watmore and paleontologist Donald Prothero reported in a preprint a giant sperm whale tooth identified as cf. Livaytan discovered in Mission Viejo, California during housing development during the 1980s and '90s. The tooth resided in the Orange County Paleontological Collection cataloged as OCPC 3125/66099 and was incomplete but nevertheless measured at least in length and in diameter. Due to poor geographic recording at the time of its discovery, the exact stratigraphic locality was unknown, but it was reported to have come from a zone that contains both the mid-Miocene Monterey Formation and younger Capistrano Formation, the latter dating between 6.6 and 5.8 mya. The authors found the preservation of the tooth to be more consistent with Capistrano Formation fossils. At the area where part of the tooth broke off revealed layers of cementum and dentin of thickness within the known range of L. melvillei teeth. OCPC 3125/66099 represented the first evidence that either Livyatan or Livyatan-like whales were not restricted to the Southern Hemisphere and likely indicated a possibly global distribution of the cetaceans.
Description
The body length of Livyatan is unknown since only the holotype skull is preserved. Lambert and colleagues estimated the body length of Livyatan using Zygophyseter and modern sperm whales as a guide. The authors opted to use the relationship between the bizygomatic width (distance between the opposite zygomatic processes) of the skull and body length because of the variable rostrum length in modern sperm whales and the rostrum of Livyatan being proportionally shorter. Doing so produced length estimates of when using the modern sperm whale and when using Zygophyseter. It has been estimated to weigh based on the length estimate of . By comparison, the modern sperm whale length measures on average for females and for males, with some males reaching up to long. The large size was probably an anti-predator adaptation, and allowed it to feed on larger prey. Livyatan is the largest fossil sperm whale discovered, and was also one of the biggest-known predators, having the largest bite of any tetrapod.
Skull
The holotype skull of Livyatan was about long. Like other raptorial sperm whales, Livyatan had a wide gap in between the temporal fossae on the sides of the skull and the zygomatic processes on the front of the skull, indicating a large space for holding strong temporal muscles, which are the most powerful muscles between the skull and the jaw. The snout was robust, thick and relatively short, which allowed it to clamp down harder and better handle struggling prey. The left and right premaxillae on the snout probably did not intersect at the tip of the snout, though the premaxillae took up most of the front end of the snout. Unlike in the modern sperm whale, the premaxillae reached the sides of the snout. The upper jaw was thick, especially midway through the snout. The snout was asymmetrical, with the right maxilla in the upper jaw becoming slightly convex towards the back of the snout, and the left maxilla becoming slightly concave towards the back of the snout. The vomer reached the tip of the snout, and was slightly concave, decreasing in thickness from the back to the front. A sudden thickening in the middle-left side of the vomer may indicate the location of the nose plug muscles. Each mandible in the lower jaw was higher than it was wide, with a larger gap in between the two than in the modern sperm whale. The mandibular symphysis which connects the two halves of the mandibles in the middle of the lower jaw was unfused. The condyloid process, which connects the lower jaw to the skull, was located near the bottom of the mandible like other sperm whales.
Teeth
Unlike the modern sperm whale, Livyatan had functional teeth in both jaws. The wearing on the teeth indicates that the teeth sheared past each other while biting down, meaning it could bite off large portions of flesh from its prey. Also, the teeth were deeply embedded into the gums and could interlock, which were adaptations to holding struggling prey. None of the teeth of the holotype were complete, and none of the back teeth were well-preserved. The lower jaw contained 22 teeth, and the upper jaw contained 18 teeth. Unlike other sperm whales with functional teeth in the upper jaw, none of the tooth roots were entirely present in the premaxilla portion of the snout, being at least partially in the maxilla. Consequently, its tooth count was lower than those sperm whales, and, aside from the modern dwarf (Kogia sima) and pygmy (K. breviceps) sperm whales, it had the lowest tooth count in the lower jaw of any sperm whale.
The most robust teeth in Livyatan were the fourth, fifth and sixth teeth in each side of the jaw. The well-preserved teeth all had a height greater than , and the largest teeth of the holotype were the second and third on the left lower jaw, which were calculated to be around high. The first right tooth was the smallest at around . The Beaumaris sperm whale tooth measured around in length, and is the largest fossil tooth discovered in Australia. These teeth are thought to be among the largest of any known animal, excluding tusks. Some of the lower teeth have been shown to contain a facet for when the jaws close, which may have been used to properly fit the largest teeth inside the jaw. In the front teeth, the tooth diameter decreased towards the base. This was the opposite for the back teeth, and the biggest diameters for these teeth were around in the lower jaw. All teeth featured a rapid shortening of the diameter towards the tip of the tooth, which were probably in part due to wearing throughout their lifetimes. The curvature of the teeth decreased from front to back, and the lower teeth were more curved at the tips than the upper teeth. The front teeth projected forward at a 45° angle, and, as in other sperm whales, cementum was probably added onto the teeth throughout the animal's lifetime.
All tooth sockets were cylindrical and single-rooted. The tooth sockets increased in size from the first to the fourth and then decreased, the fourth being the largest at around in diameter in the upper jaws, which is the largest of any known whale species. The tooth sockets were smaller in the lower jaw than they were in the upper jaw, and they were circular in shape, except for the front sockets which were more ovular.
Basin
The fossil skull of Livyatan had a curved basin, known as the supracranial basin, which was deep and wide. Unlike other raptorial sperm whales, but much like in the modern sperm whale, the basin spanned the entire length of the snout, causing the entire skull to be concave on the top rather than creating a snout as seen in Zygophyseter and Acrophyseter. The supracranial basin was the deepest and widest over the braincase, and, unlike other raptorial sperm whales, it did not overhang the eye socket. It was defined by high walls on the sides. The antorbital notches, which are usually slit-like notches on the sides of the skull right before the snout, were inside the basin. A slanting crest on the temporal fossa directed towards the back of the skull separated the snout from the rest of the skull, and was defined by a groove starting at the antorbital processes on the cheekbones. The basin had two foramina in the front, as opposed to the modern sperm whale which has one foramen on the maxilla, and to the modern dwarf and pygmy sperm whales which have several in the basin. The suture in the basin between the maxilla and the forehead had an interlocking pattern.
Classification
Livyatan was part of a fossil stem group of hyper-predatory sperm whales commonly known as macroraptorial sperm whales, or raptorial sperm whales, alongside the extinct whales Brygmophyseter, Acrophyseter and Zygophyseter. This group is known for having large, functional teeth in both the upper and lower jaws, which were used in capturing large prey, and had an enamel coating. Conversely, the modern sperm whale (Physeter macrocephalus) lacks teeth in the upper jaw, and the ability to use its teeth to catch prey. Livyatan belongs to a different lineage in respect to the other raptorial sperm whales, and the size increase and the development of the spermaceti organ, an organ that is characteristic of sperm whales, are thought to have evolved independently from other raptorial sperm whales. The large teeth of the raptorial sperm whales either evolved once in the group with a basilosaurid-like common ancestor, or independently in Livyatan. The large temporal fossa in the skull of raptorial sperm whales is thought to a plesiomorphic feature, that is, a trait inherited from a common ancestor. Since the teeth of foetal modern sperm whales (Physeter macrocephalus) have enamel on them before being coated with cementum, it is thought that the enamel is also an ancient characteristic (basal). The appearance of raptorial sperm whales in the fossil record coincides with the diversification of baleen whales in the Miocene, implying that they evolved specifically to exploit baleen whales. It has also been suggested that the raptorial sperm whales should be placed into the subfamily Hoplocetinae, alongside the genera Diaphorocetus, Idiorophus, Scaldicetus and Hoplocetus, which are known from the Miocene to the lower Pliocene. However, most of these taxa remain too fragmentary or have been used as wastebasket taxa for non-diagnostic material of stem physeteroids. This subfamily is characterized by their robust and enamel-coated teeth.
The cladogram below is modified from Lambert et al. (2017) and represents the phylogenetic relationships between Livyatan and other sperm whales, with genera identified as macroraptorial sperm whales in bold.
Palaeobiology
Hunting
Livyatan was an apex predator, and probably had a profound impact on the structuring of Miocene marine communities. Using its large and deeply rooted teeth, it is likely to have hunted large prey near the surface, its diet probably consisting mainly of medium-sized baleen whales ranging from in length. It probably also preyed upon sharks, seals, dolphins and other large marine vertebrates, occupying a niche similar to the modern killer whale (Orcinus orca). It was contemporaneous with and occupied the same region as the otodontid shark O. megalodon, which was likely also an apex predator, implying competition over their similar food sources. It is assumed that the hunting tactics of Livyatan for hunting whales were similar to that of the modern killer whale, pursuing prey to wear it out, and then drowning it. Modern killer whales work in groups to isolate and kill whales, but, given its size, Livyatan may have been able to hunt alone.
Isotopic analysis of enamel from a tooth from Chile revealed that this individual likely operated at latitudes south of 40°S. Isotopic analyses of contemporary baleen whales in the same formation show that this Livyatan was not commonly feeding on them, indicating it probably did not exclusively eat large prey, though it may have targeted baleen whales from higher latitudes.
Spermaceti organ
The supracranial basin in its head suggests that Livyatan had a large spermaceti organ, a series of oil and wax reservoirs separated by connective tissue. The uses for the spermaceti organ in Livyatan are unknown. Much like in the modern sperm whale, it could have been used in the process of biosonar to generate sound for locating prey. It is possible that it was also used as a means of acoustic displays, such as for communication purposes between individuals. It may have been used for acoustic stunning, which would have caused the bodily functions of a target animal to shut down from exposure to the intense sounds.
Another theory says that the enlarged forehead caused by the presence of the spermaceti organ is used in all sperm whales between males fighting for females during mating season by head-butting each other, including Livyatan and the modern sperm whale. It may have also been used to ram into prey; if this is the case, in support of this, there have been two reports of modern sperm whales attacking whaling vessels by ramming into them, and the organ is disproportionally larger in male modern sperm whales.
An alternate theory is that sperm whales, including Livyatan, can alter the temperature of the wax in the organ to aid in buoyancy. Lowering the temperature increases the density to have it act as a weight for deep-sea diving, and raising the temperature decreases the density to have it pull the whale to the surface.
Palaeoecology
Fossils conclusively identified as L. melvillei have been found in Peru and Chile. However, additional isolated large sperm whale teeth from other locations including California, Australia, Argentina and South Africa have been identified as a species or possible close relative of Livyatan. On the basis of these fossils, it was likely that the distribution of Livyatan was widespread. Prior to 2023, paleontologists initially believed that the genus was restricted to the Southern Hemisphere. The warmer waters around the equator have been known to be a climatic barrier for numerous cetaceans since Neogene times, and it was then-hypothesized is that Livyatan may have been among the cetaceans unable to cross the equatorial barrier. However, collecting bias was another explanation given the apparent rarity and poor fossil record of Livyatan, now supported by the Northern Hemisphere occurrence in California.
The holotype of L. melvillei is from the Tortonian stage of the Upper Miocene 9.9–8.9 mya in the Pisco Formation of Peru, which is known for its well-preserved assemblage of marine vertebrates. Among the baleen whales found, the most common was an undescribed species of cetotheriid whale measuring around , and most of the other baleen whales found were roughly the same size. Toothed whale remains found consist of beaked whales (such as Messapicetus gregarius), ancient pontoporiids (such as Brachydelphis mazeasi), oceanic dolphins and the raptorial sperm whale Acrophyseter. All seal remains found represent the earless seals. Also found were large sea turtles such as Pacifichelys urbinai, which points to the development of seagrasses in this area. Partial bones of crocodiles were discovered. Of the seabirds, fragmentary bones of cormorants and petrels were discovered, as well as two species of boobies. The remains of many cartilaginous fish were discovered in this formation, including more than 3,500 shark teeth, which mainly belonged to the ground sharks, such as requiem sharks and hammerhead sharks. To a lesser extent, mackerel sharks were also found, such as white sharks, sand sharks and Otodontidae. Many shark teeth were associated with the extinct broad-tooth mako (Cosmopolitodus/Carcharodon hastalis) and megalodon, and the teeth of these two sharks were found near whale and seal remains. Eagle rays, sawfish and angelsharks were other cartilaginous fish found. Most of the bony fish findings belonged to tunas and croakers. Livyatan and megalodon were likely the apex predators of this area during this time.
L. melvillei is also known from the Bahía Inglesa Formation of Chile, whose fossiliferous beds are dated between the Tortonian and Messinian 9.03–6.45 mya. Like the Pisco Formation, the Bahía Inglesa Formation famously holds one of the richest marine vertebrate assemblages. Baleen whale remains include ancient minke whales, grey whales, bowhead whales and cetotheriids. Of the toothed whales, five species of pontoporiids as well as beaked whales, porpoises, three other species of sperm whales such as cf. Scaldicetus, and the Odobenocetops have been yielded. Other marine mammals include the marine sloth Thalassocnus and pinnipeds like Acrophoca. At least 28 different species of sharks have been described, including many extant ground sharks and white sharks as well as extinct species such as the false mako (Parotodus sp.), broad-toothed mako, megalodon and the transitional great white Carcharodon hubbelli. Other marine vertebrates include penguins and other seabirds, and species of crocodiles and ghavials.
The Beaumaris sperm whale was found in the Beaumaris Bay Black Rock Sandstone Formation in Australia near the city of Melbourne, dating to 5mya in the Pliocene. Beaumaris Bay is one of the most productive marine fossil sites in Australia for marine megafauna. Shark teeth belonging to twenty different species have been discovered there, such as from the whale shark (Rhincodon typus), the Port Jackson shark (Heterodontus portusjacksoni), the broad-toothed mako and megalodon. Some examples of whales found include the ancient humpback whale Megaptera miocaena, the dolphin Steno cudmorei and the sperm whale Physetodon baileyi. Other large marine animals found include ancient elephant seals, dugongs, sea turtles, ancient penguins such as Pseudaptenodytes, the extinct albatross Diomedea thyridata and the extinct toothed seabirds of the genus Pelagornis.
The South African teeth attributed as cf. Livyatan are from the Avontuur Member of the Alexander Bay Formation near the village of Hondeklip Bay, Namaqualand, which is also dated to around 5mya in the Pliocene. The Hondeklip Bay locality enjoys a rich heritage of marine fossils, whose diversity may have been thanks to the initiation of the Benguela Upwelling during the late Miocene, which likely provided large populations of phytoplankton traveling the cold nutrient-rich waters. Cetaceans are the most abundant fauna in the bay, although remains tend to be difficult to conclusively identify. Included are three species of balaenopterids including two undetermined species and one identified as cf. Plesiobalaenoptera, an ancient grey whale (cf. Eschrichtius sp.), an undetermined balaenid, an unidentified dolphin, and another undetermined species of macroraptorial sperm whale. Other localities of similar age on the South African west coast have also yielded many additional species of balaenopterids and sperm whales as well as ten species of beaked whales. Large sperm whale teeth of up to around ~ in length are common in Hondeklip Bay, indicating a high presence of large sperm whales like Livyatan in the area. The locality has also a high presence of sharks indicated by a large abundance of shark teeth; however, most of these teeth have not been identified. Megalodon teeth have been found in the bay, and evidence from bite marks in whale bones indicate the additional presence of the great white shark, shortfin mako and broad-toothed mako. Other marine fauna known in Hondeklip Bay include pinnipeds such as Homiphoca capensis, bony fish and rays.
Extinction
Livyatan-like sperm whales became extinct by the early Pliocene likely due to a cooling trend causing baleen whales to increase in size and decrease in diversity, becoming coextinct with the smaller whales they fed on. Their extinction also coincides with the emergence of the orcas as well as large predatory globicephaline dolphins, possibly acting as an additional stressor to their already collapsing niche.
| Biology and health sciences | Cetaceans | Animals |
5308894 | https://en.wikipedia.org/wiki/Space%20%28mathematics%29 | Space (mathematics) | In mathematics, a space is a set (sometimes known as a universe) endowed with a structure defining the relationships among the elements of the set.
A subspace is a subset of the parent space which retains the same structure.
While modern mathematics uses many types of spaces, such as Euclidean spaces, linear spaces, topological spaces, Hilbert spaces, or probability spaces, it does not define the notion of "space" itself.
A space consists of selected mathematical objects that are treated as points, and selected relationships between these points. The nature of the points can vary widely: for example, the points can represent numbers, functions on another space, or subspaces of another space. It is the relationships that define the nature of the space. More precisely, isomorphic spaces are considered identical, where an isomorphism between two spaces is a one-to-one correspondence between their points that preserves the relationships. For example, the relationships between the points of a three-dimensional Euclidean space are uniquely determined by Euclid's axioms, and all three-dimensional Euclidean spaces are considered identical.
Topological notions such as continuity have natural definitions for every Euclidean space. However, topology does not distinguish straight lines from curved lines, and the relation between Euclidean and topological spaces is thus "forgetful". Relations of this kind are treated in more detail in the "Types of spaces" section.
It is not always clear whether a given mathematical object should be considered as a geometric "space", or an algebraic "structure". A general definition of "structure", proposed by Bourbaki, embraces all common types of spaces, provides a general definition of isomorphism, and justifies the transfer of properties between isomorphic structures.
History
Before the golden age of geometry
In ancient Greek mathematics, "space" was a geometric abstraction of the three-dimensional reality observed in everyday life. About 300 BC, Euclid gave axioms for the properties of space. Euclid built all of mathematics on these geometric foundations, going so far as to define numbers by comparing the lengths of line segments to the length of a chosen reference segment.
The method of coordinates (analytic geometry) was adopted by René Descartes in 1637. At that time, geometric theorems were treated as absolute objective truths knowable through intuition and reason, similar to objects of natural science; and axioms were treated as obvious implications of definitions.
Two equivalence relations between geometric figures were used: congruence and similarity. Translations, rotations and reflections transform a figure into congruent figures; homotheties — into similar figures. For example, all circles are mutually similar, but ellipses are not similar to circles. A third equivalence relation, introduced by Gaspard Monge in 1795, occurs in projective geometry: not only ellipses, but also parabolas and hyperbolas, turn into circles under appropriate projective transformations; they all are projectively equivalent figures.
The relation between the two geometries, Euclidean and projective, shows that mathematical objects are not given to us with their structure. Rather, each mathematical theory describes its objects by some of their properties, precisely those that are put as axioms at the foundations of the theory.
Distances and angles cannot appear in theorems of projective geometry, since these notions are neither mentioned in the axioms of projective geometry nor defined from the notions mentioned there. The question "what is the sum of the three angles of a triangle" is meaningful in Euclidean geometry but meaningless in projective geometry.
A different situation appeared in the 19th century: in some geometries the sum of the three angles of a triangle is well-defined but different from the classical value (180 degrees). Non-Euclidean hyperbolic geometry, introduced by Nikolai Lobachevsky in 1829 and János Bolyai in 1832 (and Carl Friedrich Gauss in 1816, unpublished) stated that the sum depends on the triangle and is always less than 180 degrees. Eugenio Beltrami in 1868 and Felix Klein in 1871 obtained Euclidean "models" of the non-Euclidean hyperbolic geometry, and thereby completely justified this theory as a logical possibility.
This discovery forced the abandonment of the pretensions to the absolute truth of Euclidean geometry. It showed that axioms are not "obvious", nor "implications of definitions". Rather, they are hypotheses. To what extent do they correspond to an experimental reality? This important physical problem no longer has anything to do with mathematics. Even if a "geometry" does not correspond to an experimental reality, its theorems remain no less "mathematical truths".
A Euclidean model of a non-Euclidean geometry is a choice of some objects existing in Euclidean space and some relations between these objects that satisfy all axioms (and therefore, all theorems) of the non-Euclidean geometry. These Euclidean objects and relations "play" the non-Euclidean geometry like contemporary actors playing an ancient performance. Actors can imitate a situation that never occurred in reality. Relations between the actors on the stage imitate relations between the characters in the play. Likewise, the chosen relations between the chosen objects of the Euclidean model imitate the non-Euclidean relations. It shows that relations between objects are essential in mathematics, while the nature of the objects is not.
The golden age of geometry and afterwards
The word "geometry" (from Ancient Greek: geo- "earth", -metron "measurement") initially meant a practical way of processing lengths, regions and volumes in the space in which we live, but was then extended widely (as well as the notion of space in question here).
According to Bourbaki, the period between 1795 (Géométrie descriptive of Monge) and 1872 (the "Erlangen programme" of Klein) can be called "the golden age of geometry". The original space investigated by Euclid is now called three-dimensional Euclidean space. Its axiomatization, started by Euclid 23 centuries ago, was reformed with Hilbert's axioms, Tarski's axioms and Birkhoff's axioms. These axiom systems describe the space via primitive notions (such as "point", "between", "congruent") constrained by a number of axioms.
Analytic geometry made great progress and succeeded in replacing theorems of classical geometry with computations via invariants of transformation groups. Since that time, new theorems of classical geometry have been of more interest to amateurs than to professional mathematicians. However, the heritage of classical geometry was not lost. According to Bourbaki, "passed over in its role as an autonomous and living science, classical geometry is thus transfigured into a universal language of contemporary mathematics".
Simultaneously, numbers began to displace geometry as the foundation of mathematics. For instance, in Richard Dedekind's 1872 essay Stetigkeit und irrationale Zahlen (Continuity and irrational numbers), he asserts that points on a line ought to have the properties of Dedekind cuts, and that therefore a line was the same thing as the set of real numbers. Dedekind is careful to note that this is an assumption that is incapable of being proven. In modern treatments, Dedekind's assertion is often taken to be the definition of a line, thereby reducing geometry to arithmetic. Three-dimensional Euclidean space is defined to be an affine space whose associated vector space of differences of its elements is equipped with an inner product. A definition "from scratch", as in Euclid, is now not often used, since it does not reveal the relation of this space to other spaces. Also, a three-dimensional projective space is now defined as the space of all one-dimensional subspaces (that is, straight lines through the origin) of a four-dimensional vector space. This shift in foundations requires a new set of axioms, and if these axioms are adopted, the classical axioms of geometry become theorems.
A space now consists of selected mathematical objects (for instance, functions on another space, or subspaces of another space, or just elements of a set) treated as points, and selected relationships between these points. Therefore, spaces are just mathematical structures of convenience. One may expect that the structures called "spaces" are perceived more geometrically than other mathematical objects, but this is not always true.
According to the famous inaugural lecture given by Bernhard Riemann in 1854, every mathematical object parametrized by n real numbers may be treated as a point of the n-dimensional space of all such objects. Contemporary mathematicians follow this idea routinely and find it extremely suggestive to use the terminology of classical geometry nearly everywhere.
Functions are important mathematical objects. Usually they form infinite-dimensional function spaces, as noted already by Riemann and elaborated in the 20th century by functional analysis.
Taxonomy of spaces
Three taxonomic ranks
While each type of space has its own definition, the general idea of "space" evades formalization. Some structures are called spaces, other are not, without a formal criterion. Moreover, there is no consensus on the general idea of "structure".
According to Pudlák, "Mathematics [...] cannot be explained completely by a single concept such as the mathematical structure. Nevertheless, Bourbaki's structuralist approach is the best that we have."
We will return to Bourbaki's structuralist approach in the last section "Spaces and structures", while we now outline a possible classification of spaces (and structures) in the spirit of Bourbaki.
We classify spaces on three levels. Given that each mathematical theory describes its objects by some of their properties, the first question to ask is: which properties? This leads to the first (upper) classification level. On the second level, one takes into account answers to especially important questions (among the questions that make sense according to the first level). On the third level of classification, one takes into account answers to all possible questions.
For example, the upper-level classification distinguishes between Euclidean and projective spaces, since the distance between two points is defined in Euclidean spaces but undefined in projective spaces.
Another example. The question "what is the sum of the three angles of a triangle" makes sense in a Euclidean space but not in a projective space. In a non-Euclidean space the question makes sense but is answered differently, which is not an upper-level distinction.
Also, the distinction between a Euclidean plane and a Euclidean 3-dimensional space is not an upper-level distinction; the question "what is the dimension" makes sense in both cases.
The second-level classification distinguishes, for example, between Euclidean and non-Euclidean spaces; between finite-dimensional and infinite-dimensional spaces; between compact and non-compact spaces, etc.
In Bourbaki's terms, the second-level classification is the classification by "species". Unlike biological taxonomy, a space may belong to several species.
The third-level classification distinguishes, for example, between spaces of different dimension, but does not distinguish between a plane of a three-dimensional Euclidean space, treated as a two-dimensional Euclidean space, and the set of all pairs of real numbers, also treated as a two-dimensional Euclidean space. Likewise it does not distinguish between different Euclidean models of the same non-Euclidean space.
More formally, the third level classifies spaces up to isomorphism. An isomorphism between two spaces is defined as a one-to-one correspondence between the points of the first space and the points of the second space, that preserves all relations stipulated according to the first level. Mutually isomorphic spaces are thought of as copies of a single space. If one of them belongs to a given species then they all do.
The notion of isomorphism sheds light on the upper-level classification. Given a one-to-one correspondence between two spaces of the same upper-level class, one may ask whether it is an isomorphism or not. This question makes no sense for two spaces of different classes.
An isomorphism to itself is called an automorphism. Automorphisms of a Euclidean space are shifts, rotations, reflections and compositions of these. Euclidean space is homogeneous in the sense that every point can be transformed into every other point by some automorphism.
Euclidean axioms leave no freedom; they determine uniquely all geometric properties of the space. More exactly: all three-dimensional Euclidean spaces are mutually isomorphic. In this sense we have "the" three-dimensional Euclidean space. In Bourbaki's terms, the corresponding theory is univalent. In contrast, topological spaces are generally non-isomorphic; their theory is multivalent. A similar idea occurs in mathematical logic: a theory is called categorical if all its models of the same cardinality are mutually isomorphic. According to Bourbaki, the study of multivalent theories is the most striking feature which distinguishes modern mathematics from classical mathematics.
Relations between species of spaces
Topological notions (continuity, convergence, open sets, closed sets etc.) are defined naturally in every Euclidean space. In other words, every Euclidean space is also a topological space. Every isomorphism between two Euclidean spaces is also an isomorphism between the corresponding topological spaces (called "homeomorphism"), but the converse is wrong: a homeomorphism may distort distances. In Bourbaki's terms, "topological space" is an underlying structure of the "Euclidean space" structure. Similar ideas occur in category theory: the category of Euclidean spaces is a concrete category over the category of topological spaces; the forgetful (or "stripping") functor maps the former category to the latter category.
A three-dimensional Euclidean space is a special case of a Euclidean space. In Bourbaki's terms, the species of three-dimensional Euclidean space is richer than the species of Euclidean space. Likewise, the species of compact topological space is richer than the species of topological space.
Such relations between species of spaces may be expressed diagrammatically as shown in Fig. 3. An arrow from A to B means that every is also a or may be treated as a or provides a etc. Treating A and B as classes of spaces one may interpret the arrow as a transition from A to B. (In Bourbaki's terms, "procedure of deduction" of a from a Not quite a function unless the classes A,B are sets; this nuance does not invalidate the following.) The two arrows on Fig. 3 are not invertible, but for different reasons.
The transition from "Euclidean" to "topological" is forgetful. Topology distinguishes continuous from discontinuous, but does not distinguish rectilinear from curvilinear. Intuition tells us that the Euclidean structure cannot be restored from the topology. A proof uses an automorphism of the topological space (that is, self-homeomorphism) that is not an automorphism of the Euclidean space (that is, not a composition of shifts, rotations and reflections). Such transformation turns the given Euclidean structure into a (isomorphic but) different Euclidean structure; both Euclidean structures correspond to a single topological structure.
In contrast, the transition from "3-dim Euclidean" to "Euclidean" is not forgetful; a Euclidean space need not be 3-dimensional, but if it happens to be 3-dimensional, it is full-fledged, no structure is lost. In other words, the latter transition is injective (one-to-one), while the former transition is not injective (many-to-one). We denote injective transitions by an arrow with a barbed tail, "↣" rather than "→".
Both transitions are not surjective, that is, not every B-space results from some A-space. First, a 3-dim Euclidean space is a special (not general) case of a Euclidean space. Second, a topology of a Euclidean space is a special case of topology (for instance, it must be non-compact, and connected, etc). We denote surjective transitions by a two-headed arrow, "↠" rather than "→". See for example Fig. 4; there, the arrow from "real linear topological" to "real linear" is two-headed, since every real linear space admits some (at least one) topology compatible with its linear structure.
Such topology is non-unique in general, but unique when the real linear space is finite-dimensional. For these spaces the transition is both injective and surjective, that is, bijective; see the arrow from "finite-dim real linear topological" to "finite-dim real linear" on Fig. 4. The inverse transition exists (and could be shown by a second, backward arrow). The two species of structures are thus equivalent. In practice, one makes no distinction between equivalent species of structures. Equivalent structures may be treated as a single structure, as shown by a large box on Fig. 4.
The transitions denoted by the arrows obey isomorphisms. That is, two isomorphic lead to two isomorphic .
The diagram on Fig. 4 is commutative. That is, all directed paths in the diagram with the same start and endpoints lead to the same result. Other diagrams below are also commutative, except for dashed arrows on Fig. 9. The arrow from "topological" to "measurable" is dashed for the reason explained there: "In order to turn a topological space into a measurable space one endows it with a σ-algebra. The σ-algebra of Borel sets is the most popular, but not the only choice." A solid arrow denotes a prevalent, so-called "canonical" transition that suggests itself naturally and is widely used, often implicitly, by default. For example, speaking about a continuous function on a Euclidean space, one need not specify its topology explicitly. In fact, alternative topologies exist and are used sometimes, for example, the fine topology; but these are always specified explicitly, since they are much less notable that the prevalent topology. A dashed arrow indicates that several transitions are in use and no one is quite prevalent.
Types of spaces
Linear and topological spaces
Two basic spaces are linear spaces (also called vector spaces) and topological spaces.
Linear spaces are of algebraic nature; there are real linear spaces (over the field of real numbers),
complex linear spaces (over the field of complex numbers), and more generally, linear spaces over any field. Every complex linear space is also a real linear space (the latter underlies the former), since each complex number can be specified by two real numbers. For example, the complex plane treated as a one-dimensional complex linear space may be downgraded to a two-dimensional real linear space. In contrast, the real line can be treated as a one-dimensional real linear space but not a complex linear space. | Mathematics | Set theory | null |
37285387 | https://en.wikipedia.org/wiki/Hipparionini | Hipparionini | Hipparionini is a tribe of three-toed horses in the subfamily Equinae. They had body forms similar to modern equines, with high-crowned teeth. They first appeared in North America during the Early Miocene around 17 million years ago, before migrating into the Old World around 11.4-11.0 million years ago. The youngest species date to the Early Pleistocene, becoming extinct following the arrives of modern equines of the genus Equus to the Old World.
Description
Hipparionines varied widely in size, with the smallest species like Hipparion periafricanum having a body mass of only , considerably smaller than living equines, while the largest species had body masses over .
Evolutionary history
In North America, hipparionins were equally diverse to equins during the Middle Miocene but overtook them in species richness during the Late Miocene and Early Pliocene. At the end of the Hemphillian, hipparionins severely declined in diversity.
Ecology
In the Old World hipparionins were initially browsers and mixed feeders (both browsing and grazing), over time there was increasing proportion of pure grazers, though the groups ecology remained diverse, with mixed feeding being the dominant ecology during the Pliocene. Hipparionins in the western Mediterranean during the Vallesian and Turolian exhibited noticeable niche partitioning, with smaller forms being mixed feeders while larger species had more grazing diets. In contrast, contemporaneous eastern Mediterranean hipparionins did not exhibit such niche partitioning.
Taxonomy
North American genera:
"Hipparion" (distinct from Old World species assigned to this genus)
Neohipparion
Pseudhipparion
Nannippus
Cormohipparion
Old World genera: (widely thought to descend from Cormohipparion)
Hipparion sensu stricto
Hippotherium
Cremohipparion
Sivalhippus
Eurygnathohippus
Plesiohipparion
Proboscidipparion
Shanxihippus
| Biology and health sciences | Equidae | Animals |
31696675 | https://en.wikipedia.org/wiki/Geomagnetic%20pole | Geomagnetic pole | The geomagnetic poles are antipodal points where the axis of a best-fitting dipole intersects the surface of Earth. This theoretical dipole is equivalent to a powerful bar magnet at the center of Earth, and comes closer than any other point dipole model to describing the magnetic field observed at Earth's surface. In contrast, the magnetic poles of the actual Earth are not antipodal; that is, the line on which they lie does not pass through Earth's center.
Owing to the motion of fluid in the Earth's outer core, the actual magnetic poles are constantly moving (secular variation). However, over thousands of years, their direction averages to the Earth's rotation axis. On the order of once every half a million years, the poles reverse (i.e., north switches place with south) although the time frame of this switching can be anywhere from every 10 thousand years to every 50 million years. The poles also swing in an oval of around in diameter daily due to solar wind deflecting the magnetic field.
Although the geomagnetic pole is only theoretical and cannot be located directly, it arguably is of more practical relevance than the magnetic (dip) pole. This is because the poles describe a great deal about the Earth's magnetic field, determining for example where auroras can be observed. The dipole model of the Earth's magnetic field consists of the location of geomagnetic poles and the dipole moment, which describes the strength of the field.
Definition
As a first-order approximation, the Earth's magnetic field can be modeled as a simple dipole (like a bar magnet), tilted about 9.6° with respect to the Earth's rotation axis (which defines the Geographic North and Geographic South Poles) and centered at the Earth's center. The North and South Geomagnetic Poles are the antipodal points where the axis of this theoretical dipole intersects the Earth's surface. Thus, unlike the actual magnetic poles, the geomagnetic poles always have an equal degree of latitude and supplementary degrees of longitude respectively (2017: Lat. 80.5°N, 80.5°S; Long. 72.8°W, 107.2°E). If the Earth's magnetic field were a perfect dipole, the field lines would be vertical to the surface at the Geomagnetic Poles, and they would align with the North and South magnetic poles, with the North Magnetic Pole at the south end of dipole. However, the approximation is imperfect, and so the Magnetic and Geomagnetic Poles lie some distance apart.
Location
Like the North Magnetic Pole, the North Geomagnetic Pole attracts the north pole of a bar magnet and so is in a physical sense actually a magnetic south pole. It is the center of the 'open' magnetic field lines which connect to the interplanetary magnetic field and provide a direct route for the solar wind to reach the ionosphere. , it was located at , on Ellesmere Island, Nunavut, Canada, compared to 2015, when it was located at , also on Ellesmere Island.
The South Geomagnetic Pole is the point where the axis of this best-fitting tilted dipole intersects the Earth's surface in the southern hemisphere. , it is located at , whereas in 2005, it was calculated to be located at , near Vostok Station.
Because the Earth's actual magnetic field is not an exact dipole, the (calculated) North and South Geomagnetic Poles do not coincide with the North and South Magnetic Poles. If the Earth's magnetic fields were exactly dipolar, the north pole of a magnetic compass needle would point directly at the North Geomagnetic Pole. In practice, it does not because the geomagnetic field that originates in the core has a more complex non-dipolar part, and magnetic anomalies in the Earth's crust also contribute to the local field.
The locations of geomagnetic poles are calculated by a statistical fit to measurements of the Earth's field by satellites and in geomagnetic observatories. This can be the International Geomagnetic Reference Field (covering a wide time-span in history) or the U.S. World Magnetic Model (only covering a five-year period).
Movement
The geomagnetic poles move over time because the geomagnetic field is produced by motion of the molten iron alloys in the Earth's outer core. (See geodynamo.) Over the past 150 years, the poles have moved westward at a rate of 0.05° to 0.1° per year and closer to the true poles at 0.01° per year.
Over several thousand years, the average location of the geomagnetic poles coincides with the geographical poles. Paleomagnetists have long relied on the geocentric axial dipole (GAD) hypothesis, which states that — aside from during geomagnetic reversals — the time-averaged position of the geomagnetic poles has always coincided with the geographic poles. There is considerable paleomagnetic evidence supporting this hypothesis.
Geomagnetic reversal
Over the life of the Earth, the orientation of Earth's magnetic field has reversed many times, with geomagnetic north becoming geomagnetic south and vice versa – an event known as a geomagnetic reversal. Evidence of geomagnetic reversals can be seen at mid-ocean ridges where tectonic plates move apart. As magma seeps out of the mantle and solidifies to become new ocean floor, the magnetic minerals in it are magnetized in the direction of the magnetic field. The study of this remanence is called palaeomagnetism. Thus, starting at the most recently formed ocean floor, one can read out the direction of the magnetic field in previous times as one moves farther away to older ocean floor.
| Physical sciences | Geophysics | Earth science |
21629856 | https://en.wikipedia.org/wiki/Manchester%20computers | Manchester computers | The Manchester computers were an innovative series of stored-program electronic computers developed during the 30-year period between 1947 and 1977 by a small team at the University of Manchester, under the leadership of Tom Kilburn. They included the world's first stored-program computer, the world's first transistorised computer, and what was the world's fastest computer at the time of its inauguration in 1962.
The project began with two aims: to prove the practicality of the Williams tube, an early form of computer memory based on standard cathode-ray tubes (CRTs); and to construct a machine that could be used to investigate how computers might be able to assist in the solution of mathematical problems. The first of the series, the Manchester Baby, ran its first program on 21 June 1948. As the world's first stored-program computer, the Baby, and the Manchester Mark 1 developed from it, quickly attracted the attention of the United Kingdom government, who contracted the electrical engineering firm of Ferranti to produce a commercial version. The resulting machine, the Ferranti Mark 1, was the world's first commercially available general-purpose computer.
The collaboration with Ferranti eventually led to an industrial partnership with the computer company ICL, who made use of many of the ideas developed at the university, particularly in the design of their 2900 series of computers during the 1970s.
Manchester Baby
The Manchester Baby was designed as a test-bed for the Williams tube, an early form of computer memory, rather than as a practical computer. Work on the machine began in 1947, and on 21 June 1948 the computer successfully ran its first program, consisting of 17 instructions written to find the highest proper factor of 218 (262,144) by trying every integer from 218 − 1 downwards. The program ran for 52 minutes before producing the correct answer of 217 (131,072).
The Baby was in length, tall, and weighed almost 1 long ton. It contained 550 thermionic valves – 300 diodes and 250 pentodes – and had a power consumption of 3.5 kilowatts. Its successful operation was reported in a letter to the journal Nature published in September 1948, establishing it as the world's first stored-program computer. It quickly evolved into a more practical machine, the Manchester Mark 1.
Manchester Mark 1
Development of the Manchester Mark 1 began in August 1948, with the initial aim of providing the university with a more realistic computing facility. In October 1948 UK Government Chief Scientist Ben Lockspeiser was given a demonstration of the prototype, and was so impressed that he immediately initiated a government contract with the local firm of Ferranti to make a commercial version of the machine, the Ferranti Mark 1.
Two versions of the Manchester Mark 1 were produced, the first of which, the Intermediary Version, was operational by April 1949. The Final Specification machine, which was fully working by October 1949, contained 4,050 valves and had a power consumption of 25 kilowatts. Perhaps the Manchester Mark 1's most significant innovation was its incorporation of index registers, commonplace on modern computers.
In June 2022 an IEEE Milestone was dedicated to the "Manchester University "Baby" Computer and its Derivatives, 1948-1951".
Meg and Mercury
As a result of experience gained from the Mark 1, the developers concluded that computers would be used more in scientific roles than pure maths. They therefore embarked on the design of a new machine which would include a floating-point unit; work began in 1951. The resulting machine, which ran its first program in May 1954, was known as Meg, or the megacycle machine. It was smaller and simpler than the Mark 1, as well as quicker at solving maths problems. Ferranti produced a commercial version marketed as the Ferranti Mercury, in which the Williams tubes were replaced by the more reliable core memory.
Transistor Computer
Work on building a smaller and cheaper computer began in 1952, in parallel with Meg's ongoing development. Two of Kilburn's team, Richard Grimsdale and D. C. Webb, were assigned to the task of designing and building a machine using the newly developed transistors instead of valves, which became known as the Manchester TC. Initially the only devices available were germanium point-contact transistors; these were less reliable than the valves they replaced but consumed far less power.
Two versions of the machine were produced. The first was the world's first transistorised computer, a prototype, and became operational on 16 November 1953. "The 48-bit machine used 92 point-contact transistors and 550 diodes". The second version was completed in April 1955. The 1955 version used 250 junction transistors, 1,300 solid-state diodes, and had a power consumption of 150 watts. The machine did however make use of valves to generate its 125 kHz clock waveforms and in the circuitry to read and write on its magnetic drum memory, so it was not the first completely transistorised computer, a distinction that went to the Harwell CADET of 1955.
Problems with the reliability of early batches of transistors meant that the machine's mean time between failures was about 90 minutes, which improved once the more reliable junction transistors became available. The Transistor Computer's design was adopted by the local engineering firm of Metropolitan-Vickers in their Metrovick 950, in which all the circuitry was modified to make use of junction transistors. Six Metrovick 950s were built, the first completed in 1956. They were successfully deployed within various departments of the company and were in use for about five years.
Muse and Atlas
Development of MUSE – a name derived from "microsecond engine" – began at the university in 1956. The aim was to build a computer that could operate at processing speeds approaching one microsecond per instruction, one million instructions per second. Mu (or μ) is a prefix in the SI and other systems of units denoting a factor of 10−6 (one millionth).
At the end of 1958 Ferranti agreed to collaborate with Manchester University on the project, and the computer was shortly afterwards renamed Atlas, with the joint venture under the control of Tom Kilburn. The first Atlas was officially commissioned on 7 December 1962, and was considered at that time to be the most powerful computer in the world, equivalent to four IBM 7094s. It was said that whenever Atlas went offline half of the UK's computer capacity was lost. Its fastest instructions took 1.59 microseconds to execute, and the machine's use of virtual storage and paging allowed each concurrent user to have up to one million words of storage space available. Atlas pioneered many hardware and software concepts still in common use today including the Atlas Supervisor, "considered by many to be the first recognisable modern operating system".
Two other machines were built: one for a joint British Petroleum/University of London consortium, and the other for the Atlas Computer Laboratory at Chilton near Oxford. A derivative system was built by Ferranti for Cambridge University, called the Titan or Atlas 2, which had a different memory organisation, and ran a time-sharing operating system developed by Cambridge Computer Laboratory.
The University of Manchester's Atlas was decommissioned in 1971, but the last was in service until 1974. Parts of the Chilton Atlas are preserved by the National Museums of Scotland in Edinburgh.
In June 2022 an IEEE Milestone was dedicated to the "Atlas Computer and the Invention of Virtual Memory 1957–1962".
MU5
The Manchester MU5 was the successor to Atlas. An outline proposal for a successor to Atlas was presented at the 1968 IFIP Conference in Edinburgh, although work on the project and talks with ICT (of which Ferranti had become part) aimed at obtaining their assistance and support had begun in 1966. The new machine, later to become known as MU5, was intended to be at the top end of a range of machines and to be 20 times faster than Atlas.
In 1968 the Science Research Council (SRC) awarded Manchester University a five-year grant of £630,466 (equivalent to £ in ) to develop the machine and ICT, later to become ICL, made its production facilities available to the University. In that year a group of 20 people was involved in the design: 11 Department of Computer Science staff, 5 seconded ICT staff and 4 SRC supported staff. The peak level of staffing was in 1971, when the numbers, including research students, rose to 60.
The most significant novel features of the MU5 processor were its instruction set and the use of associative memory to speed up operand and instruction accesses. The instruction set was designed to permit the generation of efficient object code by compilers, to allow for a pipeline organisation of the processor and to provide information to the hardware on the nature of operands, so as to allow them to be optimally buffered. Thus named variables were buffered separately from array elements, which were themselves accessed by means of named descriptors. Each descriptor included an array length which could be used in string processing instructions or to enable array bound checking to be carried out by hardware. The instruction pre-fetching mechanism used an associative jump trace to predict the outcome of impending branches.
The MU5 operating system MUSS was designed to be highly adaptable and was ported to a variety of processors at Manchester and elsewhere. In the completed MU5 system, three processors (MU5 itself, an ICL 1905E and a PDP-11), as well as a number of memories and other devices, were interconnected by a high-speed Exchange. All three processors ran a version of MUSS. MUSS also encompassed compilers for various languages and runtime packages to support the compiled code. It was structured as a small kernel that implemented an arbitrary set of virtual machines analogous to a corresponding set of processors. The MUSS code appeared in the common segments that formed part of each virtual machine's virtual address space.
MU5 was fully operational by October 1974, coinciding with ICL's announcement that it was working on the development of a new range of computers, the 2900 series. ICL's 2980 in particular, first delivered in June 1975, owed a great deal to the design of MU5. MU5 remained in operation at the University until 1982. A fuller article about MU5 can be found on the Engineering and Technology History Wiki.
MU6
Once MU5 was fully operational, a new project was initiated to produce its successor, MU6. MU6 was intended to be a range of processors: MU6P, an advanced microprocessor architecture intended for use as a personal computer,
MU6-G, a high performance machine for general or scientific applications and MU6V, a parallel vector processing system. A prototype model of MU6V, based on 68000 microprocessors with vector orders emulated as "extracodes" was constructed and tested but not further developed beyond this. MU6-G was built with a grant from SRC and successfully ran as a service machine in the Department between 1982 and 1987, using the MUSS operating system developed as part of the MU5 project.
SpiNNaker
SpiNNaker: Spiking Neural Network Architecture is a massively parallel, manycore supercomputer architecture designed by Steve Furber in the University of Manchester's Advanced Processor Technologies Research Group (APT). Built in 2019, it is composed of 57,600 ARM9 processors (specifically ARM968), each with 18 cores and 128 MB of mobile DDR SDRAM, totalling 1,036,800 cores and over 7 TB of RAM. The computing platform is based on spiking neural networks, useful in simulating the human brain (see Human Brain Project).
Summary
| Technology | Early computers | null |
26010282 | https://en.wikipedia.org/wiki/Brane | Brane | In string theory and related theories (such as supergravity theories), a brane is a physical object that generalizes the notion of a zero-dimensional point particle, a one-dimensional string, or a two-dimensional membrane to higher-dimensional objects. Branes are dynamical objects which can propagate through spacetime according to the rules of quantum mechanics. They have mass and can have other attributes such as charge.
Mathematically, branes can be represented within categories, and are studied in pure mathematics for insight into homological mirror symmetry and noncommutative geometry.
The word "brane" originated in 1987 as a contraction of "membrane".
p-branes
A point particle is a 0-brane, of dimension zero; a string, named after vibrating musical strings, is a 1-brane; a membrane, named after vibrating membranes such as drumheads, is a 2-brane. The corresponding object of arbitrary dimension p is called a p-brane, a term coined by M. J. Duff et al. in 1988.
A p-brane sweeps out a (p+1)-dimensional volume in spacetime called its worldvolume. Physicists often study fields analogous to the electromagnetic field, which live on the worldvolume of a brane.
D-branes
In string theory, a string may be open (forming a segment with two endpoints) or closed (forming a closed loop). D-branes are an important class of branes that arise when one considers open strings. As an open string propagates through spacetime, its endpoints are required to lie on a D-brane. The letter "D" in D-brane refers to the Dirichlet boundary condition, which the D-brane satisfies.
One crucial point about D-branes is that the dynamics on the D-brane worldvolume is described by a gauge theory, a kind of highly symmetric physical theory which is also used to describe the behavior of elementary particles in the standard model of particle physics. This connection has led to important insights into gauge theory and quantum field theory. For example, it led to the discovery of the AdS/CFT correspondence, a theoretical tool that physicists use to translate difficult problems in gauge theory into more mathematically tractable problems in string theory.
Categorical description
Mathematically, branes can be described using the notion of a category. This is a mathematical structure consisting of objects, and for any pair of objects, a set of morphisms between them. In most examples, the objects are mathematical structures (such as sets, vector spaces, or topological spaces) and the morphisms are functions between these structures. One can likewise consider categories where the objects are D-branes and the morphisms between two branes and are states of open strings stretched between and .
In one version of string theory known as the topological B-model, the D-branes are complex submanifolds of certain six-dimensional shapes called Calabi–Yau manifolds, together with additional data that arise physically from having charges at the endpoints of strings. Intuitively, one can think of a submanifold as a surface embedded inside of a Calabi–Yau manifold, although submanifolds can also exist in dimensions different from two. In mathematical language, the category having these branes as its objects is known as the derived category of coherent sheaves on the Calabi–Yau. In another version of string theory called the topological A-model, the D-branes can again be viewed as submanifolds of a Calabi–Yau manifold. Roughly speaking, they are what mathematicians call special Lagrangian submanifolds. This means, among other things, that they have half the dimension of the space in which they sit, and they are length-, area-, or volume-minimizing. The category having these branes as its objects is called the Fukaya category.
The derived category of coherent sheaves is constructed using tools from complex geometry, a branch of mathematics that describes geometric shapes in algebraic terms and solves geometric problems using algebraic equations. On the other hand, the Fukaya category is constructed using symplectic geometry, a branch of mathematics that arose from studies of classical physics. Symplectic geometry studies spaces equipped with a symplectic form, a mathematical tool that can be used to compute area in two-dimensional examples.
The homological mirror symmetry conjecture of Maxim Kontsevich states that the derived category of coherent sheaves on one Calabi–Yau manifold is equivalent in a certain sense to the Fukaya category of a completely different Calabi–Yau manifold. This equivalence provides an unexpected bridge between two branches of geometry, namely complex and symplectic geometry.
| Physical sciences | Particle physics: General | Physics |
40078184 | https://en.wikipedia.org/wiki/Tinder%20%28app%29 | Tinder (app) | Tinder is an online dating and geosocial networking application launched in 2012. On Tinder, users "swipe right" to like or "swipe left" to dislike other users' profiles, which include their photos, a short bio, and some of their interests. Tinder uses a "double opt-in" system, also called "matching", where two users must like each other before they can exchange messages.
In 2022, Tinder had 10.9 million subscribers and 75 million monthly active users. As of 2021, Tinder had recorded more than 65 billion matches worldwide.
History
Sean Rad launched Tinder at a hackathon at the Hatch Labs incubator in West Hollywood, Los Angeles County, California, United States in 2012.
Sean Rad and engineer Joe Munoz built the original prototype for Tinder, "MatchBox", during a hackathon in February 2012. The hackathon was hosted by Hatch Labs, a New York-based startup incubator with a West Hollywood outpost. Realizing the name MatchBox was too similar to Match.com, Rad, his co-founders, and early employees renamed the company Tinder. The company's flame-themed logo remained consistent throughout the rebranding.
2012: Prototype and launch
In January 2012, Rad was hired as general manager of Cardify, a credit card loyalty app launched by Hatch Labs. During a hackathon in his first month, he presented the idea for a dating app called Matchbox. Rad and Munoz built the prototype for MatchBox and presented the app on February 16, 2012.
In March, co-founder Jonathan Badeen (front-end operator and later Tinder's CSO), and Chris Gulczynski (designer and later Tinder's CCO) joined Cardify.
In May, while Cardify was going through Apple's App Store approval process, the team focused on MatchBox. During the same period, Alexa Mateen (Justin's sister) and her friend Whitney Wolfe Herd were hired as Cardify sales reps.
In August 2012, Cardify was abandoned, Matchbox was renamed Tinder, and co-founder Justin Mateen (marketer and later Tinder's CMO) joined the company.
In September 2012, Tinder was soft-launched in the App Store. It was then launched at several college campuses and started to expand quickly.
2013: Swipe feature developed
Tinder's selection function, which was initially click-based, evolved into the company's swipe feature. The feature was established when Rad and Badeen, interested in gamification, modeled the feature on a deck of cards. Badeen then streamlined the action after a trial on a bathroom mirror. Tinder has been credited with popularizing the swipe feature many companies now use.
2014–2016: Growth
By October 2014, Tinder users completed over one billion swipes per day, producing about 12 million matches per day. By then, Tinder's average user generally spent about 90 minutes a day on the app.
Rad served as Tinder's CEO until March 2015, when he was replaced by former eBay and Microsoft executive Chris Payne. Rad returned as CEO in August 2015.
In 2015, Tinder released its "Rewind" function and its "Super Like" function and retired its Tinder "Moments" and "Last Active" feature. In January 2015, Tinder acquired Chill, the developer of Tappy, a mobile messenger that uses "images and ephemerality".
In 2016, Tinder was the most popular dating app in the United States, holding a 25.6% market share of monthly users. On the company's third-quarter earnings call, Match Group's CEO Greg Blatt described the popular dating app Tinder as a "rocket" and the "future of this business."
In September 2016, the company also initiated testing of its "Boost" functionality in Australia. The feature went live for all users in October of that year.
In October 2016, Tinder announced the opening of its first office in Silicon Valley in the hope of more effectively recruiting technical employees.
In November 2016, Tinder introduced more options for gender selection.
In December 2016, Blatt took over as interim CEO of Tinder. Rad stepped down as CEO, becoming chairman of the company.
2017: Merger with Match
Tinder had annual revenue of $403 million and accounted for 31% of Match Group's 2017 annual revenue of $1.28 billion. In the same year, it surpassed Netflix as the highest-grossing app on the app store. Match Group's market cap as of December 28, 2017, was $10.03 billion.
In March 2017, Tinder launched Tinder Online, a web-optimized version of the app. Initially, it was available only in Argentina, Brazil, Colombia, Indonesia, Italy, Mexico, the Philippines, and Sweden, and did not include special features such as "Super Likes" or "Tinder Boost". Tinder Online launched globally in September 2017. During the launch of the web version, Tinder took legal action to shut down third-party apps providing a web extension to use the Tinder app from a desktop computer.
In July 2017, Match Group merged with Tinder for approximately $3 billion.
In August 2017, Tinder launched an additional subscription service, Tinder Gold, that allows subscribers to see which users "swiped back" without alerting those users. Tinder Gold quickly became a successful revenue source for the company, boosting Match Group's total revenue by 19% compared to 2016. This boost in revenue and profits came as Tinder's paid member count rose by a record 476,000 to more than 2.5 million, mainly due to product changes and technology improvements. Tinder Gold's popularity led to a surge in Match Group shares and recorded high share prices. Blatt called Tinder's performance "fantastic" and said the company was driving most of Match Group's growth in late 2017.
Blatt resigned from Match Group and Tinder in 2017 after allegations of sexual harassment. He was replaced by Elie Seidman.
2018–2019
In 2018, Tinder had annual revenue of $805 million and accounted for 48% of Match Group's 2018 annual revenue of $1.67 billion. Match Group's market cap as of December 30, 2018, was $15.33 billion.
On August 6, 2018, Tinder had over 3.7 million paid subscribers, up 81% over the same quarter in 2017. On August 21, 2018, Tinder launched Tinder University, a feature that allows college students to connect with other students on their campus and at nearby schools.
In 2019, Tinder had annual revenue of $1.152 billion and accounted for 58% of Match Group's total 2019 annual revenue of $2.0 billion. Match Group's market cap as of December 30, 2019, was $21.09 billion.
On May 10, it was reported that Tinder was planning for a lighter version of the app, Tinder Lite, aimed at growing markets where data usage, bandwidth and storage space are a concern.
Tinder had 5.2 million paying subscribers at the end of 2019's second quarter, up 1.5 million from the year-ago quarter and up 503,000 from the first quarter of 2019. It became the highest-grossing non-gaming app, beating Netflix. Tinder's subscriber growth led Match Group's shares to the best single-day gain in their history on August 7, adding more than $5 billion to the company's market capitalization.
On September 12, Tinder relaunched Swipe Night, an interactive series where users make decisions following a storyline. Swipe Night had previously been launched in October 2019. It was slated to be launched internationally in March 2020, but it was postponed until September due to the COVID-19 pandemic. Swipe Nights international launch included multiple countries and languages.
In 2019, Tinder had annual revenue of $1.2 billion.
2020
In 2020, Tinder had annual revenue of $1.355 billion and accounted for 58% of Match Group's 2020 revenue of $2.34 billion. Match Group's market cap as of December 23, 2020, was $40.45 billion.
In January 2020, the Tinder administration enabled a panic button and anti-catfishing technology to improve U.S. users' safety. In the future, these features are planned to become globally available. If something goes wrong on a date, a user can hit a panic button, transmit accurate location data, and call emergency services. To use this feature, users must download and install the Noonlight app. Also, before going to a meeting, users are required to take selfies to prove their photos in Tinder profiles match their real identities.
In response to the global COVID-19 pandemic, in March 2020 Tinder temporarily made its Passport feature available for free to all its users worldwide. Previously this feature had been only accessible to users who had purchased a subscription.
In August, Tinder revealed plans for its Platinum subscription plan, which gives users access to more features for a higher price than gold. The same month, Jim Lanzone took over as CEO.
On September 1, Tinder was banned in Pakistan in a crackdown on what the Pakistani government deemed "immoral content".
On November 4, Tinder reported higher than expected third-quarter earnings and significant platform growth amid the pandemic: the app grew its user base by 15% and its subscriber count by 16% since the third quarter of 2019. According to Business Insider, Tinder's growth was fueled by a large population that turned to online dating in response to increasing social isolation and health risks. In 2020, Match Group products, including Tinder, released video-based features to facilitate long-distance dating. Tinder had 6.6 million subscribers globally in November, growing from 6.2 million reported in June.
2021 to date
Match Group's market cap as of October 14, 2021, was $44.59 billion.
In February 2021, Tinder announced it would launch a range of mobile accessories under the brand name Tinder Made. The app reported that month an all-time high in users ready to "go on a date" as opposed to virtual and online chats during the height of the pandemic in the U.S. It gave away pairs of COVID-19 testing kits to some matches to encourage responsible behavior as users begin to meet in person again.
In March 2021, Tinder announced a service that would let users run background checks on potential matches after an investment in Garbo, a company that "collects public records and reports of violence or abuse, including arrests, convictions, restraining orders, harassment, and other violent crimes". Garbo does not publicize drug possession charges or traffic violations, citing disproportionate incarceration. This service comes with a fee that has not yet been disclosed to users.
In August 2021, Tinder announced an ID verification service to mitigate catfishing on the platform.
In September 2021, Lanzone announced that he was stepping down as chief executive to pursue a new role at Yahoo. Tinder named Renate Nyborg its new CEO. She is the company's first female CEO.
In December 2021, Nyborg announced that the company was working on creating a metaverse called Tinderverse, a shared virtual reality. The company is also testing Tinder Coins, an in-app currency users can earn as a reward for good behavior, allowing them to pay for the platform's premium services.
In May 2023, Match Group announced its intent to restrict access to Tinder in Russia and withdraw from the Russian market by June 30, 2023, citing the need to protect human rights. In doing so, it became one of many Western companies to leave Russia after the Russian invasion of Ukraine. Tinder became unavailable in Russia on the intended date. In response, a group of Russians staged a mock funeral in Sochi. Mourners, dressed in black, shared their stories and experiences with the app and placed red carnations on a mock gravestone in the form of a smartphone.
In 2023, Tinder introduced a new feature, Tinder Matchmaker. It allows the user's friends and family to access their Tinder account and suggest potential partners for them.
According to Tinder, the addition of a feature that allows friends to select compatible matches makes online dating safer because the "team-based" approach helps users mitigate potential risks associated with interacting with strangers. Melissa Hobley, Tinder's Chief Marketing Officer, said that singles often seek advice from friends when choosing a match: "The new feature streamlines this process, allowing your most trusted friends to join you on your dating journey."
In January 2024, Match Group named Faye Iosotaluno as chief executive officer of Tinder.
Tinder ceased operations in Belarus after February 14, 2024.
Operation
After building a profile with a Meta login or cell phone number, users can swipe yes (right) or no (left) on potential romantic matches. Chatting on Tinder is only available between two users who have swiped right on one another's profiles. The selections a user makes are not known to other users unless two users swipe right on each other's profiles. Once the user has matches, they can send personal photos, called "Tinder Moments", to all matches at once, allowing each match to like the photos. The site also has verified profiles for public figures, so that celebrities and other public figures can verify their identities.
The app is available in about 190 countries and can be used in 45 languages.
Financials
Since merging with Tinder in July 2017, Match Group's market capitalization has grown from $8.34 billion to $44.59 billion as of October 14, 2021. In August 2021, Morgan Stanley valued Tinder's worth at $42 billion. The valuation is based on a multiple of 40x EBITDA, similar to its counterpart Bumble.
In March 2014, media and internet conglomerate IAC increased its majority stake in Tinder, a move that is believed to have valued Tinder at several billion dollars. In July 2015, Tinder was valued at $1.35 billion by Bank of America Merrill Lynch based upon an estimate of $27 per user on an estimated user base of 50 million with an additional bullish estimate of $3 billion by taking the average of the IPOs of similar companies. Analysts also estimated that Tinder had about half a million paid users within its user base, which consisted mostly of free users. The monetization of the site has come through leaving the basic app free and then adding different in-app purchase options for additional functions and features.
Paid subscriptions
Tinder launched its subscription version, Tinder Plus, in March 2015, with the functionality enabling limitless matches, whereas the free Tinder app restricts the number of right swipes in a 12-hour period. It has sparked debate by restricting the number of "likes" a free user may offer in a given length of time, as well as charging varying amounts for different age groups. The price of a Tinder Plus subscription was £14.99/US$19.99 per month for users over 28 and £3.99/US$9.99 per month for users 28 and under.
In 2016, Tinder independently increased paying members by nearly one million, while Match Group's 44 other brands added just 1.4 million.
In June 2017, Tinder launched Tinder Gold, a members-only service that offers users Tinder's most exclusive features: Passport, Rewind, Unlimited Likes, Likes You, five Super Likes per day, one Boost per month, and more profile controls. The price of a monthly membership started at $14.99 for users under 30 and $29.99 for users over 30.
As of December 2020, Tinder had 6.6 million paid users. According to a Match Group SEC filing, the growth in international and North American average subscribers was primarily driven by Tinder.
Users
Tinder is used widely throughout the world and available in over 190 countries and 56 languages. As of September 2021, an estimated 75 million people used the app every month. In late 2014, Tinder users averaged 12 million matches per day. To get those 12 million matches, users collectively made around 1 billion swipes per day. Tinder now limits users' number of available swipes per 12 hours based on an algorithm to make sure users are actually looking at profiles and not just spamming the app to rack up random matches. As of June 2016, Tinder may not be used by anyone under 18. If minors are found to be using Tinder, they are banned until they are 18. As of April 2015, Tinder users swiped through 1.6 billion Tinder profiles and made more than 26 million matches per day. More than 558 billion matches have been made since Tinder launched.
Advertising
An ad campaign launched by "The Barn" internship program of Bartle Bogle Hegarty (BBH) used Tinder profiles to promote their NYC Puppy Rescue Project. BBH added Facebook pet profiles to the Tinder network. The campaign received media coverage from Slate, Inc., The Huffington Post, and others. In April 2015, Tinder revealed its first sponsored ad promoting Budweiser's next #Whatever, USA campaign.
On December 11, 2020, Tinder announced a partnership with popular artist Megan Thee Stallion for the Put Yourself Out There Challenge, giving $10,000 to users who made unique profiles.
Lawsuits
On June 30, 2014, former vice president of marketing Whitney Wolfe filed a sexual harassment and sex discrimination suit in Los Angeles County Superior Court against IAC-owned Match Group, Tinder's parent company. The lawsuit alleged that Rad and Mateen had engaged in discrimination, sexual harassment, and retaliation against her, while Tinder's corporate supervisor, IAC's Sam Yagan, did nothing. IAC suspended Mateen from his position pending an ongoing investigation, and stated that it "acknowledges that Mateen sent private messages containing inappropriate content, but it believes Mateen, Rad and the company are innocent of the allegations". The suit was settled with no admission of wrongdoing, and Wolfe reportedly received over $1 million in the settlement.
In March 2018, Match Group sued Bumble, arguing that it was guilty of patent infringement and of stealing trade secrets from Tinder. In June 2020, an undisclosed settlement was reached between Match Group and Bumble to settle all litigations.
In December 2018, The Verge reported that Tinder had dismissed its vice president of marketing and communication, Rosette Pambakian. Pambakian alleged former Match Group and IAC CEO Greg Blatt sexually assaulted her in a hotel room after a company party in December 2016. She further accused the company of firing her when she reported the incident.
In August 2018, Rad, Mateen, and eight other former and current Tinder executives filed suit against Match Group and IAC, alleging that they manipulated the 2017 valuation of the company to deny them billions of dollars they were owed. The suit charges that Match Group and IAC executives deliberately manipulated the data given to the banks, overestimating expenses and underestimating potential revenue growth, to keep the 2017 valuation artificially low. Tinder's 2017 valuation was set at $3 billion, unchanged from a valuation that had been done two years earlier, despite rapid growth in revenue and subscribers. The plaintiffs sought more than $2 billion in damages. The trial was scheduled to begin on November 8, 2021. On December 1, 2021, Match Group and the plaintiffs settled for $441 million.
Criticism
Privacy concerns
Critics have raised concerns about Tinder regarding issues including cybersecurity, data privacy, and public health. Public health officials in Rhode Island and Utah have claimed that Tinder and similar apps are responsible for an increase in some STDs. In February 2014, security researchers in New York found a flaw that made it possible to find users' precise locations for between 40 and 165 days. Tinder's spokesperson, Rosette Pambakian, said the issue was resolved within 48 hours. Tinder CEO Sean Rad said in a statement that shortly after being contacted, Tinder implemented specific measures to enhance location security and further obscure location data. In August 2016, two engineers found another flaw that showed all users' matches' exact locations. The location was updated every time a user logged into the app, even for blocked matches. The issue was detected in March 2016 and fixed in August 2016.
In July 2017, a study published in Advances in Intelligent Systems and Computing found that Tinder users are excessively willing to disclose their personally identifiable information. In September 2017, The Guardian published an article by a journalist who requested all data that the Tinder app had recorded about her from the company and found that Tinder stores all user messages, user locations and times, the characteristics of users who interest a particular user, the characteristics of particular users of interest to other users, and the length of time users spend looking at particular pictures, which for the journalist amounted to 800 pages of detail.
Safety
In 2021, Tinder partnered with and invested in Garbo, a nonprofit background check company, to add a feature enabling users to run background checks on their matches. Critics believe the integration of background check software discriminates against the third of the adult U.S. working population who have criminal records. Another issue critics raised was the background checks' unreliability, since they disproportionately impact Black people and other ethnic minorities. A Prison Policy Initiative spokesperson claimed that because the U.S. applies laws unequally, introducing criminal background checks to dating apps filters out marginalized groups. Moreover, public records and court documents often contain erroneous or outdated information. Garbo does not advertise drug possession charges or traffic violations in an attempt to combat further marginalization.
In 2022, Tinder announced a partnership with the campaign group No More in an attempt to protect its users against domestic violence, especially women, who are more vulnerable. The No More feature will educate users about safe dating.
Banning of trans people
In August 2015, Business Insider reported that transgender Tinder users were being reported and banned for being transgender. The article included an interview with a trans woman who also described abusive messages she received that included transmisogyny and homophobia. In December 2017, Vice reported that the pattern of being reported and banned had continued. The article included an anecdote from American YouTuber Kat Blaque, saying that every account she had ever had on Tinder had been banned.
In March 2018, an article in The Cut reported that a trans woman sued Tinder for removing her profile and refusing to explain why it had been deleted. The article further reported that many transgender people had their accounts reported and banned, some within several hours of opening them. In October 2019, PinkNews reported that the reporting and banning of trans people had continued. The article stated that Tinder has "50 gender options" but bans trans users for their gender identity. A trans woman is reported as saying, "the fact that the system can be abused in such a way just shows, yet again, that they [Tinder] don't care about the trans and non-binary people using their app."
In late 2019, articles in Reuters and The Independent focused on Tinder's lack of action to correct the issue of transgender users' accounts being reported and banned. In August 2022, an article in The Cut highlighted the issue again. One trans woman interviewed recommended OkCupid as friendlier to LGBTQ+ people.
Reception
Reviews
The New York Times wrote that the wide use of Tinder could be attributed not to what Tinder was doing right but to flaws in the models of earlier dating software, which relied on mathematical algorithms to select potential partners. Relationship experts interviewed by the newspaper said that users used the photos that come in succession on the app to derive cues as to social status, confidence levels, and personal interests. Marie Claire wrote that the app was "easy to use on the run," "natural," and "addictive" due to its game style, but that "it's hard to focus" and Tinder "is still very casual sex-focused—many are only on Tinder for a quick hook-up, so if it's a serious relationship you're after this app might not be for you."
In September 2020, the Pakistani government announced that it would ban five dating apps, including Tinder, because the apps provided immoral or indecent content that does not comply with Pakistani law.
User behavior
Men use dating apps and websites more than women do—both by frequency of use and number of users. According to Statista, as of March 2021, 75.8% of the U.S. Tinder user base is male and 24.2% is female. The first study on swiping strategies revealed that "men tend to like a large proportion of the women they view but receive only a tiny fraction of matches in return—just 0.6 percent", while women are much more selective about swiping but match at a 10% higher rate than men. The study also found that women are more engaged, take longer to compose their messages, and write longer messages.
According to University of Texas at Austin psychologist David Buss, "Apps like Tinder and OkCupid give people the impression that there are thousands or millions of potential mates out there. One dimension of this is the impact it has on men's psychology. When there is ... a perceived surplus of women, the whole mating system tends to shift towards short-term dating," and there is a feeling of disconnect when choosing future partners. The appearance of an abundance of potential partners may cause online daters to be less likely to choose a partner and be less satisfied with their choices of partners. Data released by Tinder in 2018 showed that of the 1.6 billion swipes it records per day, 26 million result in matches (a match rate of approximately 1.63%).
In August 2015, journalist Nancy Jo Sales wrote in Vanity Fair that Tinder operates within a culture of users seeking sex without relationships. In 2017, the Department of Communications Studies at Texas Tech University conducted a study to see how infidelity was connected to Tinder. The experiment was conducted on 550 students from an unnamed university in the southwestern United States. The students first provided their demographic information and then answered questions about Tinder's link to infidelity. The results showed that more than half reported having seen somebody on Tinder they knew was in an exclusive relationship (63.9%), while 73.1% of participants reported that they knew male friends who used Tinder while in a relationship, and 56.1% reported that they had female friends who used Tinder while in a relationship. Psychologists Douglas T. Kenrick, Sara E. Gutierres, Laurie L. Goldberg, Steven Neuberg, Kristin L. Zierk, and Jacquelyn M. Krones have demonstrated experimentally that after exposure to photos of or stories about desirable potential mates, people decrease their ratings of commitment to their partners. David Buss has estimated that approximately 30% of men on Tinder are married.
Before Tinder, most online dating services matched people according to their autobiographical information, such as interests, hobbies, and future plans. Tinder places a higher weight on first impressions. For social scientists studying human courtship behavior, Tinder offers a much simpler environment than its predecessors. In a 2016 study analyzing Tinder users' behavior in New York City and London, researchers created with three profiles using stock photographs, two with actual photographs of volunteers, one with no photos, and one that was apparently deactivated. All pictures were of people of average physical attractiveness. These profiles then liked all profiles presented to them, and then counted the number of returning likes. The researchers found that men liked a large proportion of the profiles they viewed, but received returning likes only 0.6% of the time, while women were much more selective but received matches 10% of the time. Men received matches at a much slower rate than women. Once they received a match, women were far more likely than men to send a message, 21% compared to 7%, but took more time to do so. Tyson and his team found that for the first two-thirds of messages from each sex, women sent them within 18 minutes of receiving a match compared to five minutes for men. Men's first messages had an average of 12 characters and were typical simple greetings; by contrast, initial messages by women averaged 122 characters.
By sending out questionnaires to frequent Tinder users, the researchers discovered that the reason men tended to like a large proportion of the women they saw was to increase their chances of a match. This led to a feedback loop in which men liked more and more of the profiles they saw, while women could afford to be more selective in liking profiles because of a greater probability of a match. The feedback loop's mathematical limit occurs when men like all profiles they see while women find a match whenever they like a profile. It was not known whether some evolutionarily stable strategy has emerged, nor has Tinder revealed such information.
Tyson and his team found that even though the men-to-women ratio of their data set was approximately one, the male profiles received 8,248 matches while the female profiles received 532, because the vast majority of the matches for both the male and female profiles came from male profiles (with 86% of the matches for the male profiles alone coming from other male profiles), leading the researchers to conclude that gay men were "far more active in liking than heterosexual women." On the other hand, the deactivated male account received all of its matches from women. The researchers were not sure why that happened.
| Technology | Social network and blogging | null |
3947316 | https://en.wikipedia.org/wiki/Non-stoichiometric%20compound | Non-stoichiometric compound | Non-stoichiometric compounds are chemical compounds, almost always solid inorganic compounds, having elemental composition whose proportions cannot be represented by a ratio of small natural numbers (i.e. an empirical formula); most often, in such materials, some small percentage of atoms are missing or too many atoms are packed into an otherwise perfect lattice work.
Contrary to earlier definitions, modern understanding of non-stoichiometric compounds view them as homogeneous, and not mixtures of stoichiometric chemical compounds. Since the solids are overall electrically neutral, the defect is compensated by a change in the charge of other atoms in the solid, either by changing their oxidation state, or by replacing them with atoms of different elements with a different charge. Many metal oxides and sulfides have non-stoichiometric examples; for example, stoichiometric iron(II) oxide, which is rare, has the formula , whereas the more common material is nonstoichiometric, with the formula . The type of equilibrium defects in non-stoichiometric compounds can vary with attendant variation in bulk properties of the material. Non-stoichiometric compounds also exhibit special electrical or chemical properties because of the defects; for example, when atoms are missing, electrons can move through the solid more rapidly. Non-stoichiometric compounds have applications in ceramic and superconductive material and in electrochemical (i.e., battery) system designs.
Occurrence
Iron oxides
Nonstoichiometry is pervasive for metal oxides, especially when the metal is not in its highest oxidation state. For example, although wüstite (ferrous oxide) has an ideal (stoichiometric) formula , the actual stoichiometry is closer to . The non-stoichiometry reflect the ease of oxidation of to effectively replacing a small portion of with two thirds their number of . Thus for every three "missing" ions, the crystal contains two ions to balance the charge. The composition of a non-stoichiometric compound usually varies in a continuous manner over a narrow range. Thus, the formula for wüstite is written as , where x is a small number (0.05 in the previous example) representing the deviation from the "ideal" formula. Nonstoichiometry is especially important in solid, three-dimensional polymers that can tolerate mistakes. To some extent, entropy drives all solids to be non-stoichiometric. But for practical purposes, the term describes materials where the non-stoichiometry is measurable, usually at least 1% of the ideal composition.
Iron sulfides
The monosulfides of the transition metals are often nonstoichiometric. Best known perhaps is nominally iron(II) sulfide (the mineral pyrrhotite) with a composition (x = 0 to 0.2). The rare stoichiometric endmember is known as the mineral troilite. Pyrrhotite is remarkable in that it has numerous polytypes, i.e. crystalline forms differing in symmetry (monoclinic or hexagonal) and composition (, , and others). These materials are always iron-deficient owing to the presence of lattice defects, namely iron vacancies. Despite those defects, the composition is usually expressed as a ratio of large numbers and the crystals symmetry is relatively high. This means the iron vacancies are not randomly scattered over the crystal, but form certain regular configurations. Those vacancies strongly affect the magnetic properties of pyrrhotite: the magnetism increases with the concentration of vacancies and is absent for the stoichiometric .
Palladium hydrides
Palladium hydride is a nonstoichiometric material of the approximate composition (0.02 < x < 0.58). This solid conducts hydrogen by virtue of the mobility of the hydrogen atoms within the solid.
Tungsten oxides
It is sometimes difficult to determine if a material is non-stoichiometric or if the formula is best represented by large numbers. The oxides of tungsten illustrate this situation. Starting from the idealized material tungsten trioxide, one can generate a series of related materials that are slightly deficient in oxygen. These oxygen-deficient species can be described as , but in fact they are stoichiometric species with large unit cells with the formulas , where n = 20, 24, 25, 40. Thus, the last species can be described with the stoichiometric formula , whereas the non-stoichiometric description implies a more random distribution of oxide vacancies.
Other cases
At high temperatures (1000 °C), titanium sulfides present a series of non-stoichiometric compounds.
The coordination polymer Prussian blue, nominally and their analogs are well known to form in non-stoichiometric proportions. The non-stoichiometric phases exhibit useful properties vis-à-vis their ability to bind caesium and thallium ions.
Applications
Oxidation catalysis
Many useful compounds are produced by the reactions of hydrocarbons with oxygen, a conversion that is catalyzed by metal oxides. The process operates via the transfer of "lattice" oxygen to the hydrocarbon substrate, a step that temporarily generates a vacancy (or defect). In a subsequent step, the missing oxygen is replenished by O2. Such catalysts rely on the ability of the metal oxide to form phases that are not stoichiometric. An analogous sequence of events describes other kinds of atom-transfer reactions including hydrogenation and hydrodesulfurization catalysed by solid catalysts. These considerations also highlight the fact that stoichiometry is determined by the interior of crystals: the surfaces of crystals often do not follow the stoichiometry of the bulk. The complex structures on surfaces are described by the term "surface reconstruction".
Ion conduction
The migration of atoms within a solid is strongly influenced by the defects associated with non-stoichiometry. These defect sites provide pathways for atoms and ions to migrate through the otherwise dense ensemble of atoms that form the crystals. Oxygen sensors and solid state batteries are two applications that rely on oxide vacancies. One example is the CeO2-based sensor in automotive exhaust systems. At low partial pressures of O2, the sensor allows the introduction of increased air to effect more thorough combustion.
Superconductivity
Many superconductors are non-stoichiometric. For example, yttrium barium copper oxide, arguably the most notable high-temperature superconductor, is a non-stoichiometric solid with the formula YxBa2Cu3O7−x. The critical temperature of the superconductor depends on the exact value of x. The stoichiometric species has x = 0, but this value can be as great as 1.
History
It was mainly through the work of Nikolai Semenovich Kurnakov and his students that Berthollet's opposition to Proust's law was shown to have merit for many solid compounds. Kurnakov divided non-stoichiometric compounds into berthollides and daltonides depending on whether their properties showed monotonic behavior with respect to composition or not. The term berthollide was accepted by IUPAC in 1960. The names come from Claude Louis Berthollet and John Dalton, respectively, who in the 19th century advocated rival theories of the composition of substances. Although Dalton "won" for the most part, it was later recognized that the law of definite proportions had important exceptions.
| Physical sciences | Alloys and ceramic compounds | Chemistry |
3947744 | https://en.wikipedia.org/wiki/Free%20surface%20effect | Free surface effect | The free surface effect is a mechanism which can cause a watercraft to become unstable and capsize.
It refers to the tendency of liquids — and of unbound aggregates of small solid objects, like seeds, gravel, or crushed ore, whose behavior approximates that of liquids — to move in response to changes in the attitude of a craft's cargo holds, decks, or liquid tanks in reaction to operator-induced motions (or sea states caused by waves and wind acting upon the craft). When referring to the free surface effect, the condition of a tank that is not full is described as a "slack tank", while a full tank is "pressed up".
Stability and equilibrium
In a normally loaded vessel any rolling from perpendicular is countered by a righting moment generated from the increased volume of water displaced by the hull on the lowered side. This assumes the center of gravity of the vessel is relatively constant. If a moving mass inside the vessel moves in the direction of the roll, this counters the righting effect by moving the center of gravity towards the lowered side. The free surface effect can become a problem in a craft with large partially full bulk cargo compartments, fuel tanks, or water tanks (especially if they span the full breadth of the ship), or from accidental flooding, such as has occurred in several accidents involving roll-on/roll-off ferries.
If a compartment or tank is either empty or full, there is no change in the craft's center of mass as it rolls from side to side (in strong winds, heavy seas, or on sharp motions or turns). However, if the compartment is only partially full, the liquid in the compartment will respond to the vessel's heave, pitch, roll, surge, sway or yaw. For example, as a vessel rolls to port, liquid will displace to the port side of a compartment, and this will move the vessel's center of mass to port. This has the effect of slowing the vessel's return to vertical.
The momentum of large volumes of moving liquids cause significant dynamic forces, which act against the righting effect. When the vessel returns to vertical the roll continues and the effect is repeated on the opposite side. In heavy seas, this can become a positive feedback loop, causing each roll to become more and more extreme, eventually overcoming the righting effect leading to a capsize. While repeated oscillations of increasing magnitude are commonly associated with the free surface effect, they are not a necessary condition. For example, in the cases of both the and , gradual buildup of water from fire-fighting caused capsizing in a single continuous roll.
Mitigation
To mitigate this hazard, cargo vessels use multiple smaller bulk compartments or liquid tanks, instead of fewer larger ones, and possibly baffles within bulk compartments or liquid tanks to minimize the free surface effects on the craft as a whole. Keeping individual bulk compartments or liquid tanks either relatively empty or full is another way to minimize the effect and its attendant problems. Hydraulic tankers use water to displace lighter oil to keep the tank full at all times. Tanks or compartments that do not straddle the vessel's centerline are somewhat less prone to destabilising oscillations. Similarly, narrow compartments (aligned bow to stern) and compartments at the extremes away from the centerline are less prone to cause instability.
Historical examples
Flooding, liquid cargo leakage, or unintended water (from precipitation, waves, or hull damage) in any compartments or on any decks of watercraft, and the resulting free surface effect are often a contributing cause of accidents, capsizes, and casualties e.g. the loss of (Wellington, New Zealand, April 1968), (Zeebrugge, Belgium, March 1987), and (Baltic Sea, September 1994). In the case of the RORO ferry al-Salam Boccaccio 98 (Red Sea, February 2006), improper fire-fighting procedures caused flooding leading directly to instability and capsize. In both the cases of the al-Salam Boccaccio 98 and , severe listing followed immediately after the ship had undergone a hard turn, causing unstable volumes of water (from collision damage in the latter) to surge from one side of the ship to the other.
Effects on land and aircraft
The free surface effect can affect any kind of craft, including watercraft (where it is most common), bulk cargo or liquid tanker semi-trailers and trucks (causing either jackknifing or roll-overs), and aircraft (especially fire-fighting water-droppers and refueling tankers where baffles mitigate but do not eliminate the effects). The term "free surface effect" implies a liquid under the influence of gravity. Slosh dynamics is the overarching field which covers both free surface effects and situations such as space vehicles, where gravity is inconsequential but inertia and momentum interact with complex fluid mechanics to cause vehicle instability.
Regulation
To reduce the effects of free surface effect potentially capsizing a ship, regulatory requirements apply to all ships internationally under the SOLAS Convention and the International Code on Intact Stability.
| Physical sciences | Fluid mechanics | Physics |
3949010 | https://en.wikipedia.org/wiki/Cyclic%20compound | Cyclic compound | A cyclic compound (or ring compound) is a term for a compound in the field of chemistry in which one or more series of atoms in the compound is connected to form a ring. Rings may vary in size from three to many atoms, and include examples where all the atoms are carbon (i.e., are carbocycles), none of the atoms are carbon (inorganic cyclic compounds), or where both carbon and non-carbon atoms are present (heterocyclic compounds with rings containing both carbon and non-carbon). Depending on the ring size, the bond order of the individual links between ring atoms, and their arrangements within the rings, carbocyclic and heterocyclic compounds may be aromatic or non-aromatic; in the latter case, they may vary from being fully saturated to having varying numbers of multiple bonds between the ring atoms. Because of the tremendous diversity allowed, in combination, by the valences of common atoms and their ability to form rings, the number of possible cyclic structures, even of small size (e.g., < 17 total atoms) numbers in the many billions.
Adding to their complexity and number, closing of atoms into rings may lock particular atoms with distinct substitution (by functional groups) such that stereochemistry and chirality of the compound results, including some manifestations that are unique to rings (e.g., configurational isomers). As well, depending on ring size, the three-dimensional shapes of particular cyclic structures – typically rings of five atoms and larger – can vary and interconvert such that conformational isomerism is displayed. Indeed, the development of this important chemical concept arose historically in reference to cyclic compounds. Finally, cyclic compounds, because of the unique shapes, reactivities, properties, and bioactivities that they engender, are the majority of all molecules involved in the biochemistry, structure, and function of living organisms, and in man-made molecules such as drugs, pesticides, etc.
Structure and classification
A cyclic compound or ring compound is a compound in which at least some its atoms are connected to form a ring. Rings vary in size from three to many tens or even hundreds of atoms. Examples of ring compounds readily include cases where:
all the atoms are carbon (i.e., are carbocycles),
none of the atoms are carbon (inorganic cyclic compounds), or where
both carbon and non-carbon atoms are present (heterocyclic compounds with rings containing both carbon and non-carbon).
Common atoms can (as a result of their valences) form varying numbers of bonds, and many common atoms readily form rings. In addition, depending on the ring size, the bond order of the individual links between ring atoms, and their arrangements within the rings, cyclic compounds may be aromatic or non-aromatic; in the case of non-aromatic cyclic compounds, they may vary from being fully saturated to having varying numbers of multiple bonds. As a consequence of the constitutional variability that is thermodynamically possible in cyclic structures, the number of possible cyclic structures, even of small size (e.g., <17 atoms) numbers in the many billions.
Moreover, the closing of atoms into rings may lock particular functional group–substituted atoms into place, resulting in stereochemistry and chirality being associated with the compound, including some manifestations that are unique to rings (e.g., configurational isomers); As well, depending on ring size, the three-dimensional shapes of particular cyclic structures — typically rings of five atoms and larger — can vary and interconvert such that conformational isomerism is displayed.
Carbocycles
The vast majority of cyclic compounds are organic, and of these, a significant and conceptually important portion are composed of rings made only of carbon atoms (i.e., they are carbocycles).
Inorganic cyclic compounds
Inorganic atoms form cyclic compounds as well. Examples include sulfur and nitrogen (e.g. heptasulfur imide , trithiazyl trichloride , tetrasulfur tetranitride ), silicon (e.g., cyclopentasilane ), phosphorus and nitrogen (e.g., hexachlorophosphazene ), phosphorus and oxygen (e.g., metaphosphates and other cyclic phosphoric acid derivatives), boron and oxygen (e.g., sodium metaborate , borax), boron and nitrogen (e.g. borazine ). When carbon in benzene is "replaced" by other elements, e.g., as in borabenzene, silabenzene, germanabenzene, stannabenzene, and phosphorine, aromaticity is retained, and so aromatic inorganic cyclic compounds are also known and well-characterized.
Heterocyclic compounds
A heterocyclic compound is a cyclic compound that has atoms of at least two different elements as members of its ring(s). Cyclic compounds that have both carbon and non-carbon atoms present are heterocyclic carbon compounds, and the name refers to inorganic cyclic compounds as well (e.g., siloxanes, which contain only silicon and oxygen in the rings, and borazines, which contain only boron and nitrogen in the rings). Hantzsch–Widman nomenclature is recommended by the IUPAC for naming heterocycles, but many common names remain in regular use.
Macrocycles
The term macrocycle is used for compounds having a rings of 8 or more atoms. Macrocycles may be fully carbocyclic (rings containing only carbon atoms, e.g. cyclooctane), heterocyclic containing both carbon and non-carbon atoms (e.g. lactones and lactams containing rings of 8 or more atoms), or non-carbon (containing only non-carbon atoms in the rings, e.g. diselenium hexasulfide). Heterocycles with carbon in the rings may have limited non-carbon atoms in their rings (e.g., in lactones and lactams whose rings are rich in carbon but have limited number of non-carbon atoms), or be rich in non-carbon atoms and displaying significant symmetry (e.g., in the case of chelating macrocycles). Macrocycles can access a number of stable conformations, with preference to reside in conformations that minimize transannular nonbonded interactions within the ring (e.g., with the chair and chair-boat being more stable than the boat-boat conformation for cyclooctane, because of the interactions depicted by the arcs shown). Medium rings (8-11 atoms) are the most strained, with between 9-13 (kcal/mol) strain energy, and analysis of factors important in the conformations of larger macrocycles can be modeled using medium ring conformations. Conformational analysis of odd-membered rings suggests they tend to reside in less symmetrical forms with smaller energy differences between stable conformations.
Nomenclature
IUPAC nomenclature has extensive rules to cover the naming of cyclic structures, both as core structures, and as substituents appended to alicyclic structures. The term macrocycle is used when a ring-containing compound has a ring of 12 or more atoms. The term polycyclic is used when more than one ring appears in a single molecule. Naphthalene is formally a polycyclic compound, but is more specifically named as a bicyclic compound. Several examples of macrocyclic and polycyclic structures are given in the final gallery below.
The atoms that are part of the ring structure are called annular atoms.
Isomerism
Stereochemistry
The closing of atoms into rings may lock particular atoms with distinct substitution by functional groups such that the result is stereochemistry and chirality of the compound, including some manifestations that are unique to rings (e.g., configurational isomers).
Conformational isomerism
Depending on ring size, the three-dimensional shapes of particular cyclic structures—typically rings of 5-atoms and larger—can vary and interconvert such that conformational isomerism is displayed. Indeed, the development of this important chemical concept arose, historically, in reference to cyclic compounds. For instance, cyclohexanes—six membered carbocycles with no double bonds, to which various substituents might be attached, see image—display an equilibrium between two conformations, the chair and the boat, as shown in the image.
The chair conformation is the favored configuration, because in this conformation, the steric strain, eclipsing strain, and angle strain that are otherwise possible are minimized. Which of the possible chair conformations predominate in cyclohexanes bearing one or more substituents depends on the substituents, and where they are located on the ring; generally, "bulky" substituents—those groups with large volumes, or groups that are otherwise repulsive in their interactions—prefer to occupy an equatorial location. An example of interactions within a molecule that would lead to steric strain, leading to a shift in equilibrium from boat to chair, is the interaction between the two methyl groups in cis-1,4-dimethylcyclohexane. In this molecule, the two methyl groups are in opposing positions of the ring (1,4-), and their cis stereochemistry projects both of these groups toward the same side of the ring. Hence, if forced into the higher energy boat form, these methyl groups are in steric contact, repel one another, and drive the equilibrium toward the chair conformation.
Aromaticity
Cyclic compounds may or may not exhibit aromaticity; benzene is an example of an aromatic cyclic compound, while cyclohexane is non-aromatic. In organic chemistry, the term aromaticity is used to describe a cyclic (ring-shaped), planar (flat) molecule that exhibits unusual stability as compared to other geometric or connective arrangements of the same set of atoms. As a result of their stability, it is very difficult to cause aromatic molecules to break apart and to react with other substances. Organic compounds that are not aromatic are classified as aliphatic compounds—they might be cyclic, but only aromatic rings have especial stability (low reactivity).
Since one of the most commonly encountered aromatic systems of compounds in organic chemistry is based on derivatives of the prototypical aromatic compound benzene (an aromatic hydrocarbon common in petroleum and its distillates), the word “aromatic” is occasionally used to refer informally to benzene derivatives, and this is how it was first defined. Nevertheless, many non-benzene aromatic compounds exist. In living organisms, for example, the most common aromatic rings are the double-ringed bases in RNA and DNA. A functional group or other substituent that is aromatic is called an aryl group.
The earliest use of the term “aromatic” was in an article by August Wilhelm Hofmann in 1855. Hofmann used the term for a class of benzene compounds, many of which do have odors (aromas), unlike pure saturated hydrocarbons. Today, there is no general relationship between aromaticity as a chemical property and the olfactory properties of such compounds (how they smell), although in 1855, before the structure of benzene or organic compounds was understood, chemists like Hofmann were beginning to understand that odiferous molecules from plants, such as terpenes, had chemical properties we recognize today are similar to unsaturated petroleum hydrocarbons like benzene.
In terms of the electronic nature of the molecule, aromaticity describes a conjugated system often made of alternating single and double bonds in a ring. This configuration allows for the electrons in the molecule's pi system to be delocalized around the ring, increasing the molecule's stability. The molecule cannot be represented by one structure, but rather a resonance hybrid of different structures, such as with the two resonance structures of benzene. These molecules cannot be found in either one of these representations, with the longer single bonds in one location and the shorter double bond in another (See Theory below). Rather, the molecule exhibits bond lengths in between those of single and double bonds. This commonly seen model of aromatic rings, namely the idea that benzene was formed from a six-membered carbon ring with alternating single and double bonds (cyclohexatriene), was developed by August Kekulé (see History section below). The model for benzene consists of two resonance forms, which corresponds to the double and single bonds superimposing to produce six one-and-a-half bonds. Benzene is a more stable molecule than would be expected without accounting for charge delocalization.
Principal uses
Because of the unique shapes, reactivities, properties, and bioactivities that they engender, cyclic compounds are the largest majority of all molecules involved in the biochemistry, structure, and function of living organisms, and in the man-made molecules (e.g., drugs, herbicides, etc.) through which man attempts to exert control over nature and biological systems.
Synthetic reactions
Important general reactions for forming rings
There are a variety of specialized reactions whose use is solely the formation of rings, and these will be discussed below. In addition to those, there are a wide variety of general organic reactions that historically have been crucial in the development, first, of understanding the concepts of ring chemistry, and second, of reliable procedures for preparing ring structures in high yield, and with defined orientation of ring substituents (i.e., defined stereochemistry). These general reactions include:
Acyloin condensation;
Anodic oxidations; and
the Dieckmann condensation as applied to ring formation.
Ring-closing reactions
In organic chemistry, a variety of synthetic procedures are particularly useful in closing carbocyclic and other rings; these are termed ring-closing reactions. Examples include:
alkyne trimerisation;
the Bergman cyclization of an enediyne;
the Diels–Alder, between a conjugated diene and a substituted alkene, and other cycloaddition reactions;
the Nazarov cyclization reaction, originally being the cyclization of a divinyl ketone;
various radical cyclizations;
ring-closing metathesis reactions, which also can be used to accomplish a specific type of polymerization;
the Ruzicka large ring synthesis, in which two carboxyl groups combine to form a carbonyl group with loss of and ;
the Wenker synthesis converting a beta amino alcohol to an aziridine
Ring-opening reactions
A variety of further synthetic procedures are particularly useful in opening carbocyclic and other rings, generally which contain a double bound or other functional group "handle" to facilitate chemistry; these are termed ring-opening reactions. Examples include:
ring opening metathesis, which can also be used to accomplish a specific type of polymerization.
Ring expansion and ring contraction reactions
Ring expansion and contraction reactions are common in organic synthesis, and are frequently encountered in pericyclic reactions. Ring expansions and contractions can involve the insertion of a functional group such as the case with Baeyer–Villiger oxidation of cyclic ketones, rearrangements of cyclic carbocycles as seen in intramolecular Diels-Alder reactions, or collapse or rearrangement of bicyclic compounds as several examples.
Examples
Simple, mono-cyclic examples
The following are examples of simple and aromatic carbocycles, inorganic cyclic compounds, and heterocycles:
Complex and polycyclic examples
The following are examples of cyclic compounds exhibiting more complex ring systems and stereochemical features:
| Physical sciences | Substance | Chemistry |
23139600 | https://en.wikipedia.org/wiki/Blue%20Java%20banana | Blue Java banana | The Blue Java (also known as the blue banana, Ice Cream banana, Vanilla Banana, Hawaiian banana, Ney Mannan, Krie, or Cenizo) is a hardy, cold-tolerant banana cultivar known for its sweet aromatic fruit, which is said to have an ice cream-like consistency and flavor reminiscent of vanilla. It is native to Southeast Asia and is a hybrid of two species of banana native to Southeast Asia—Musa balbisiana and Musa acuminata.
Taxonomy and nomenclature
The Blue Java banana is a triploid (ABB) hybrid of the seeded banana Musa balbisiana and Musa acuminata.
Its accepted name is Musa acuminata × balbisiana (ABB Group) 'Blue Java'.
Synonyms include:
Musa acuminata × balbisiana (ABB Group) 'Ice Cream'
In Hawaii it is known as the 'Ice Cream banana' and in Fiji as the 'Hawaiian banana'. It is also called 'Krie' in the Philippines, 'Kepok Awu' in Indonesia and 'Cenizo' in Central America.
Description
Blue Java banana trees can grow to a height of . They are cold-tolerant and wind-resistant because of their strong pseudostems and root systems. The leaves are silvery green in color.
The fruit bunches are small, bearing seven to nine hands. The fruit are in length and exhibit a characteristic silvery green color when unripe where the silveriness is caused by heavy coat of wax. The fruit turn a pale yellow when ripe, with white creamy flesh. They bloom around 15 to 24 months after planting and can be harvested after 115 to 150 days. The bananas have bumps called "knuckles" due to their passing resemblance to human knuckles.
Uses
Blue Java bananas are popular bananas that can be eaten fresh or cooked. They are known for their fragrant flavor which has a vanilla-like custard taste.
They are also popular as ornamentals and shade plants for their unusual blue coloration, large size, and tolerance to temperate climates.
Pests and diseases
Common pests
Borers
Grasshoppers
Root-knot nematode
Common diseases
Panama disease
Black sigatoka
| Biology and health sciences | Tropical and tropical-like fruit | Plants |
37299388 | https://en.wikipedia.org/wiki/Rhinocerotoidea | Rhinocerotoidea | Rhinocerotoidea is a superfamily of perissodactyls that appeared 56 million years ago in the Paleocene. They included four extinct families, the Amynodontidae, the Hyracodontidae, the Paraceratheriidae, and the Eggysodontidae. The only extant family is the Rhinocerotidae (true rhinoceroses), which survives as five living species. Extinct non-rhinocerotid members of the group are sometimes considered rhinoceroses in a broad sense. Although the term 'rhinoceroses' is sometimes used to refer to all of these, a less ambiguous vernacular term for this group is 'rhinocerotoids'. The family Paraceratheriidae contains the largest land mammals known to have ever existed.
Taxonomy
The cladogram below follows a phylogenetic analysis by Bai et al. (2020):
| Biology and health sciences | Perissodactyla | Animals |
41529802 | https://en.wikipedia.org/wiki/Giant%20Void | Giant Void | The Giant Void (also known as the Giant Void in NGH, Canes Venatici Supervoid, and AR-Lp 36) is an extremely large region of space with an underdensity of galaxies and located in the constellation Canes Venatici. It is the second-largest-confirmed void to date, with an estimated diameter of 300 to 400 Mpc (1 to 1.3 billion light-years) and its centre is approximately 1.5 billion light-years away (z = 0.116). It was discovered in 1988, and was the largest void in the Northern Galactic Hemisphere, and possibly the second-largest ever detected. Even the hypothesized "Eridanus Supervoid" corresponding to the location of the WMAP cold spot is dwarfed by this void, although the Giant Void does not correspond to any significant cooling to the cosmic microwave background.
Inside this vast void there are 17 galaxy clusters, concentrated in a spherically shaped region 50 Mpc in diameter. Studies of the motion of these clusters show that they have no interaction to each other, meaning the density of the clusters is very low resulting in weak gravitational interaction. The void's location in the sky is close to the Boötes Void.
In a series of papers published between 2004 and 2006, cosmologist and theoretical physicist Laura Mersini-Houghton presented a theory that the universe arose from a multiverse, and made a series of testable predictions which included the existence of the Giant Void.
| Physical sciences | Notable patches of universe | Astronomy |
2119174 | https://en.wikipedia.org/wiki/Effects%20of%20climate%20change | Effects of climate change | Effects of climate change are well documented and growing for Earth's natural environment and human societies. Changes to the climate system include an overall warming trend, changes to precipitation patterns, and more extreme weather. As the climate changes it impacts the natural environment with effects such as more intense forest fires, thawing permafrost, and desertification. These changes impact ecosystems and societies, and can become irreversible once tipping points are crossed. Climate activists are engaged in a range of activities around the world that seek to ameliorate these issues or prevent them from happening.
The effects of climate change vary in timing and location. Up until now the Arctic has warmed faster than most other regions due to climate change feedbacks. Surface air temperatures over land have also increased at about twice the rate they do over the ocean, causing intense heat waves. These temperatures would stabilize if greenhouse gas emissions were brought under control. Ice sheets and oceans absorb the vast majority of excess heat in the atmosphere, delaying effects there but causing them to accelerate and then continue after surface temperatures stabilize. Sea level rise is a particular long term concern as a result. The effects of ocean warming also include marine heatwaves, ocean stratification, deoxygenation, and changes to ocean currents. The ocean is also acidifying as it absorbs carbon dioxide from the atmosphere.
The ecosystems most immediately threatened by climate change are in the mountains, coral reefs, and the Arctic. Excess heat is causing environmental changes in those locations that exceed the ability of animals to adapt. Species are escaping heat by migrating towards the poles and to higher ground when they can. Sea level rise threatens coastal wetlands with flooding. Decreases in soil moisture in certain locations can cause desertification and damage ecosystems like the Amazon Rainforest. At of warming, around 10% of species on land would become critically endangered.
Humans are vulnerable to climate change in many ways. Sources of food and fresh water can be threatened by environmental changes. Human health can be impacted by weather extremes or by ripple effects like the spread of infectious diseases. Economic impacts include changes to agriculture, fisheries, and forestry. Higher temperatures will increasingly prevent outdoor labor in tropical latitudes due to heat stress. Island nations and coastal cities may be inundated by rising sea levels. Some groups of people may be particularly at risk from climate change, such as the poor, children, and indigenous peoples. Industrialised countries, which have emitted the vast majority of CO2, have more resources to adapt to global warming than developing nations do. Cumulative effects and extreme weather events can lead to displacement and migration.
Changes in temperature
Global warming affects all parts of Earth's climate system. Global surface temperatures have risen by . Scientists say they will rise further in the future. The changes in climate are not uniform across the Earth. In particular, most land areas have warmed faster than most ocean areas. The Arctic is warming faster than most other regions. Night-time temperatures have increased faster than daytime temperatures. The impact on nature and people depends on how much more the Earth warms.
Scientists use several methods to predict the effects of human-caused climate change. One is to investigate past natural changes in climate. To assess changes in Earth's past climate scientists have studied tree rings, ice cores, corals, and ocean and lake sediments. These show that recent temperatures have surpassed anything in the last 2,000 years. By the end of the 21st century, temperatures may increase to a level last seen in the mid-Pliocene. This was around 3 million years ago. At that time, mean global temperatures were about warmer than pre-industrial temperatures. The global mean sea level was up to higher than it is today. The modern observed rise in temperature and concentrations has been rapid. Even abrupt geophysical events in Earth's history do not approach current rates.
How much the world warms depends on human greenhouse gas emissions and on how sensitive the climate is to greenhouse gases. The more carbon dioxide () is emitted in the 21st century the hotter the world will be by 2100. For a doubling of greenhouse gas concentrations, the global mean temperature would rise by about . If emissions of stopped abruptly and there was no use of negative emission technologies, the Earth's climate would not start moving back to its pre-industrial state. Temperatures would stay at the same high level for several centuries. After about a thousand years, 20% to 30% of human-emitted would remain in the atmosphere. The ocean and land would not have taken them. This would commit the climate to a warmer state long after emissions have stopped.
With current mitigation policies the temperature will be about 2.7 °C (2.0–3.6 °C) above pre-industrial levels by 2100. It would rise by if governments achieved all their unconditional pledges and targets. If all the countries that have set or are considering net-zero targets achieve them, the temperature will rise by around . There is a big gap between national plans and commitments and the actions that governments have taken around the world.
Weather
The lower and middle atmosphere, where nearly all weather occurs, are heating due to the greenhouse effect. Evaporation and atmospheric moisture content increase as temperatures rise. Water vapour is a greenhouse gas, so this process is a self-reinforcing feedback.
The excess water vapour also gets caught up in storms. This makes them more intense, larger, and potentially longer-lasting. This in turn causes rain and snow events to become stronger and leads to increased risk of flooding. Extra drying worsens natural dry spells and droughts. This increases risk of heat waves and wildfires. Scientists have identified human activities as the cause of recent climate trends. They are now able to estimate the impact of climate change on extreme weather events using a process called extreme event attribution. For instance such research can look at historical data for a region and conclude that a specific heat wave was more intense due to climate change. In addition , the time shifts of the season onsets, changes in the length of the season durations have been reported in many regions of the world. As a result of changes in climatic patterns and rising global temperatures, extreme weather events like heatwaves and heavy precipitation are occurring more frequently and with increasing severity.
Heat waves and temperature extremes
Heatwaves over land have become more frequent and more intense in almost all world regions since the 1950s, due to climate change. Heat waves are more likely to occur simultaneously with droughts. Marine heatwaves are twice as likely as they were in 1980. Climate change will lead to more very hot days and fewer very cold days. There are fewer cold waves.
Experts can often attribute the intensity of individual heat waves to global warming. Some extreme events would have been nearly impossible without human influence on the climate system. A heatwave that would occur once every ten years before global warming started now occurs 2.8 times as often. Under further warming, heatwaves are set to become more frequent. An event that would occur every ten years would occur every other year if global warming reaches .
Heat stress is related to temperature. It also increases if humidity is higher. The wet-bulb temperature measures both temperature and humidity. Humans cannot adapt to a wet-bulb temperature above . This heat stress can kill people. If global warming is kept below , it will probably be possible to avoid this deadly heat and humidity in most of the tropics. But there may still be negative health impacts.
There is some evidence climate change is leading to a weakening of the polar vortex. This would make the jet stream more wavy. This would lead to outbursts of very cold winter weather across parts of Eurasia and North America and incursions of very warm air into the Arctic. Some stadies found a weakening of the AMOC by about 15% since 1950, causing cooling in the North Atlantic and warming in the Gulf Stream region. Climate change is expected to weaken AMOC in all emissions scenarios and, in some high emissions scenarios, can bring it to collapse. This can result in cooling of some parts of Europe by up to 30 degrees and warming in the southern hemisphere.
Rain
Warming increases global average precipitation. Precipitation is when water vapour condenses out of clouds, such as rain and snow. Higher temperatures increase evaporation and surface drying. As the air warms it can hold more water. For every degree Celsius it can hold 7% more water vapour. Scientists have observed changes in the amount, intensity, frequency, and type of precipitation. Overall, climate change is causing longer hot dry spells, broken by more intense rainfall.
Climate change has increased contrasts in rainfall amounts between wet and dry seasons. Wet seasons are getting wetter and dry seasons are getting drier. In the northern high latitudes, warming has also caused an increase in the amount of snow and rain. In the Southern Hemisphere, the rain associated with the storm tracks has shifted south. Changes in monsoons vary a lot. More monsoon systems are becoming wetter than drier. In Asia summer monsoons are getting wetter. The West African monsoon is getting wetter over the central Sahel, and drier in the far western Sahel.
Extreme storms
Storms become wetter under climate change. These include tropical cyclones and extratropical cyclones. Both the maximum and mean rainfall rates increase. This more extreme rainfall is also true for thunderstorms in some regions. Furthermore, tropical cyclones and storm tracks are moving towards the poles. This means some regions will see large changes in maximum wind speeds. Scientists expect there will be fewer tropical cyclones, but they expect their strength to increase. There has probably been an increase in the number of tropical cyclones that intensify rapidly. Meteorological and seismological data indicate a widespread increase in wind-driven global ocean wave energy in recent decades that has been attributed to an increase in storm intensity over the oceans due to climate change. Atmospheric turbulence dangerous for aviation (hard to predict or that cannot be avoided by flying higher) probably increases due to climate change.
Land
Floods
Due to an increase in heavy rainfall events, floods are likely to become more severe when they do occur. The interactions between rainfall and flooding are complex. There are some regions in which flooding is expected to become rarer. This depends on several factors. These include changes in rain and snowmelt, but also soil moisture. Climate change leaves soils drier in some areas, so they may absorb rainfall more quickly. This leads to less flooding. Dry soils can also become harder. In this case heavy rainfall runs off into rivers and lakes. This increases risks of flooding.
Droughts
Climate change affects many factors associated with droughts. These include how much rain falls and how fast the rain evaporates again. Warming over land increases the severity and frequency of droughts around much of the world. In some tropical and subtropical regions of the world, there will probably be less rain due to global warming. This will make them more prone to drought. Droughts are set to worsen in many regions of the world. These include Central America, the Amazon and south-western South America. They also include West and Southern Africa. The Mediterranean and south-western Australia are also some of these regions.
Higher temperatures increase evaporation. This dries the soil and increases plant stress. Agriculture suffers as a result. This means even regions where overall rainfall is expected to remain relatively stable will experience these impacts. These regions include central and northern Europe. Without climate change mitigation, around one third of land areas are likely to experience moderate or more severe drought by 2100. Due to global warming droughts are more frequent and intense than in the past.
Several impacts make their impacts worse. These are increased water demand, population growth and urban expansion in many areas. Land restoration can help reduce the impact of droughts. One example of this is agroforestry.
Wildfires
Climate change promotes the type of weather that makes wildfires more likely. In some areas, an increase of wildfires has been attributed directly to climate change. Evidence from Earth's past also shows more fire in warmer periods. Climate change increases evapotranspiration. This can cause vegetation and soils to dry out. When a fire starts in an area with very dry vegetation, it can spread rapidly. Higher temperatures can also lengthen the fire season. This is the time of year in which severe wildfires are most likely, particularly in regions where snow is disappearing.
Weather conditions are raising the risks of wildfires. In the United States, about 3 million acres burned each year from 1985 to 1995, but since 2004, between 4 and 10 million acres have burned each year, with the trend continuing to increase. However, the total area burnt by wildfires has decreased worldwide. This is mostly because savanna has been converted to cropland, so there are fewer trees to burn. Prescribed burning is an indigenous practice in the US and Australia which can reduce wildfire burning.
The carbon released from wildfires adds to carbon dioxide in Earth's atmosphere and therefore contributes to the greenhouse effect. Climate models do not yet fully reflect this climate change feedback.
Seismic and volcanic activity
In regions sensitive to climate change the frequency and intensity of eruptions will change as global warming increases. Glacier retreat and stronger precipitation can increase the chances for an eruption. As of 2024, government agencies are already addressing these changes and scientists are working to map the volcanoes most sensitive to climate change. The concerns regions are where glaciers are melting fast, and there are volcanoes heavily affected by precipitation. 716 volcanoes worldwide, may be affected by more extreme precipitation. Melting ice and extreme rainfall also increase secondary hazards, particularly lahars and disturb eruption forecasting by inducing ground displacements.
Earthquakes can be triggered by changes in the amount of stress on a fault in the Earth's crust. Strong rain, snow, drought and more pumping of groundwater by humans during droughts, can do it by increasing or reducing the weight of water on some pieces of the Earth's crust. So, as climate change will cause more extreme weather, it can induce more earthquakes. Glacier retreat reduce stress loads on Earth’s crust underneath, creating glacial earthquakes. Glacial earthquakes in Greenland for example, peak in frequency in the summer months and are increasing over time, possibly in response to global warming.
Sea level rise can also create pressure on tectonic faults, increasing risk for earthquekes.
In Greenland, the climate crisis triggered a landslide, which caused a mega-tsunami in September 2023. The event was so huge that the entire planet vibrated for nine days. Earthquake sensors around the world detected the Earth's vibration but the planetary-scale of the event was so unprecedented that at first scientists failed to understand it. Further investigation revealed that the culprit was collapse of a 1,200-metre-high mountain peak into the remote Dickson Fjord on September 16th, 2023 after the glacier below the mountain melted to a sufficient degree. The collapse into the fjord, in turn, launched a wave 200 metres high, which caused repeated movement of water back and forth in the fjord, blasting seismic waves planet wide for nine days.
Oceans
Sea level rise
Ice and snow
The cryosphere, the area of the Earth covered by snow or ice, is extremely sensitive to changes in global climate. There has been an extensive loss of snow on land since 1981. Some of the largest declines have been observed in the spring. During the 21st century, snow cover is projected to continue its retreat in almost all regions.
Glaciers decline
Since the beginning of the twentieth century, there has been a widespread retreat of glaciers. Those glaciers that are not associated with the polar ice sheets lost around 8% of their mass between 1971 and 2019. In the Andes in South America and in the Himalayas in Asia, the retreat of glaciers could impact water supply. The melting of those glaciers could also cause landslides or glacial lake outburst floods.
Ice sheets decline
The melting of the Greenland and West Antarctic ice sheets will continue to contribute to sea level rise over long time-scales. The Greenland ice sheet loss is mainly driven by melt from the top. Antarctic ice loss is driven by warm ocean water melting the outlet glaciers.
Future melt of the West Antarctic ice sheet is potentially abrupt under a high emission scenario, as a consequence of a partial collapse. Part of the ice sheet is grounded on bedrock below sea level. This makes it possibly vulnerable to the self-enhancing process of marine ice sheet instability. Marine ice cliff instability could also contribute to a partial collapse. But there is limited evidence for its importance. A partial collapse of the ice sheet would lead to rapid sea level rise and a local decrease in ocean salinity. It would be irreversible for decades and possibly even millennia. The complete loss of the West Antarctic ice sheet would cause over of sea level rise.
In contrast to the West Antarctic ice sheet, melt of the Greenland ice sheet is projected to take place more gradually over millennia. Sustained warming between (low confidence) and (medium confidence) would lead to a complete loss of the ice sheet. This would contribute to sea levels globally. The ice loss could become irreversible due to a further self-enhancing feedback. This is called the elevation-surface mass balance feedback. When ice melts on top of the ice sheet, the elevation drops. Air temperature is higher at lower altitudes, so this promotes further melting.
Sea ice decline
Sea ice reflects 50% to 70% of the incoming solar radiation back into space. Only 6% of incoming solar energy is reflected by the ocean. As the climate warms, the area covered by snow or sea ice decreases. After sea ice melts, more energy is absorbed by the ocean, so it warms up. This ice-albedo feedback is a self-reinforcing feedback of climate change. Large-scale measurements of sea ice have only been possible since satellites came into use.
Sea ice in the Arctic has declined in recent decades in area and volume due to climate change. It has been melting more in summer than it refreezes in winter. The decline of sea ice in the Arctic has been accelerating during the early twenty-first century. It has a rate of decline of 4.7% per decade. It has declined over 50% since the first satellite records. Ice-free summers are expected to be rare at degrees of warming. They are set to occur at least once every decade with a warming level of . The Arctic will likely become ice-free at the end of some summers before 2050.
Sea ice extent in Antarctica varies a lot year by year. This makes it difficult to determine a trend, and record highs and record lows have been observed between 2013 and 2023. The general trend since 1979, the start of the satellite measurements, has been roughly flat. Between 2015 and 2023, there has been a decline in sea ice, but due to the high variability, this does not correspond to a significant trend.
Permafrost thawing
Globally, permafrost warmed by about between 2007 and 2016. The extent of permafrost has been falling for decades. More decline is expected in the future. Permafrost thaw makes the ground weaker and unstable. The thaw can seriously damage human infrastructure in permafrost areas such as railways, settlements and pipelines. Thawing soil can also release methane and from decomposing microbes. This can generate a strong feedback loop to global warming. Some scientists believe that carbon storage in permafrost globally is approximately 1600 gigatons. This is twice the atmospheric pool.
Wildlife and nature
Recent warming has had a big effect on natural biological systems. Species worldwide are moving poleward to colder areas. On land, species may move to higher elevations. Marine species find colder water at greater depths. Climate change had the third biggest impact on nature out of various factors in the five decades up to 2020. Only change in land use and sea use and direct exploitation of organisms had a bigger impact.
The impacts of climate change on nature are likely to become bigger in the next few decades. The stresses caused by climate change, combine with other stresses on ecological systems such as land conversion, land degradation, harvesting, and pollution. They threaten substantial damage to unique ecosystems. They can even result in their complete loss and the extinction of species. This can disrupt key interactions between species within ecosystems. This is because species from one location do not leave the warming habitat at the same rate. The result is rapid changes in the way the ecosystem functions. Impacts include changes in regional rainfall patterns. Another is earlier leafing of trees and plants over many regions. Movements of species to higher latitudes and altitudes, changes in bird migrations, and shifting of the oceans' plankton and fish from cold- to warm-adapted communities are other impacts.
These changes of land and ocean ecosystems have direct effects on human well-being. For instance, ocean ecosystems help with coastal protection and provide food. Freshwater and land ecosystems can provide water for human consumption. Furthermore, these ecosystems can store carbon. This helps to stabilize the climate system.
Ecosystems on land
Climate change is a major driver of biodiversity loss in different land types. These include cool conifer forests, savannas, mediterranean-climate systems, tropical forests, and the Arctic tundra. In other ecosystems, land-use change may be a stronger driver of biodiversity loss, at least in the near term. Beyond 2050, climate change may be the major cause of biodiversity loss globally. Climate change interacts with other pressures. These include habitat modification, pollution and invasive species. Through this interaction, climate change increases the risk of extinction for many terrestrial and freshwater species. At of warming (around 2023) some ecosystems are threatened by mass die-offs of trees and from heatwaves. At of warming, around 10% of species on land would become critically endangered. This differs by group. For instance insects and salamanders are more vulnerable.
Rainfall on the Amazon rainforest is recycled when it evaporates back into the atmosphere instead of running off away from the rainforest. This water is essential for sustaining the rainforest. Due to deforestation the rainforest is losing this ability. This effect is even worse because climate change brings more frequent droughts to the area. The higher frequency of droughts in the first two decades of the 21st century and other data signal that a tipping point from rainforest to savanna might be close. A 2019 study concluded that this ecosystem could begin a 50-year-long collapse to a savanna around 2021. After that it would become increasingly and disproportionally more difficult to prevent or reverse this shift.
Marine ecosystems
Marine heatwaves are happening more often. They have widespread impacts on life in the oceans. These include mass dying events and coral bleaching. Harmful algae blooms have increased. This is in response to warming waters, loss of oxygen and eutrophication. Melting sea ice destroys habitat, including for algae that grows on its underside.
Ocean acidification can harm marine organisms in various ways. Shell-forming organisms like oysters are particularly vulnerable. Some phytoplankton and seagrass species may benefit. However, some of these are toxic to fish phytoplankton species. Their spread poses risks to fisheries and aquaculture. Fighting pollution can reduce the impact of acidification.
Warm-water coral reefs are very sensitive to global warming and ocean acidification. Coral reefs provide a habitat for thousands of species. They provide ecosystem services such as coastal protection and food. But 70–90% of today's warm-water coral reefs will disappear even if warming is kept to . Coral reefs are framework organisms. They build physical structures that form habitats for other sea creatures. Other framework organisms are also at risk from climate change. Mangroves and seagrass are considered to be at moderate risk from lower levels of global warming.
Tipping points and irreversible impacts
The climate system exhibits "threshold behavior" or tipping points when parts of the natural environment enter into a new state. Examples are the runaway loss of ice sheets or the dieback of forests. Tipping behavior is found in all parts of the climate system. These include ecosystems, ice sheets, and the circulation of the ocean and atmosphere. Tipping points are studied using data from Earth's distant past and by physical modeling. There is already moderate risk of global tipping points at above pre-industrial temperatures. That becomes a high risk at . It is possible that some tipping points are close or have already been crossed. Examples are the West Antarctic and Greenland ice sheets, the Amazon rainforest, and warm-water coral reefs.
Tipping points are perhaps the most dangerous aspect of future climate change, potentially leading to irreversible impacts on society. A collapse of the Atlantic meridional overturning circulation would likely halve rainfall in India and lead to severe drops in temperature in Northern Europe. Many tipping points are interlinked such that triggering one may lead to a cascade of effects. This remains a possibility even well below of warming. A 2018 study states that 45% of environmental problems, including those caused by climate change, are interconnected. This increases the risk of a domino effect.
Further impacts may be irreversible, at least over the timescale of many human generations. This includes warming of the deep ocean and acidification. These are set to continue even when global temperatures stop rising. In biological systems, the extinction of species would be an irreversible impact. In social systems, unique cultures may be lost. Climate change could make it more likely that endangered languages disappear.
Health, food security and water security
Humans have a climate niche. This is a certain range of temperatures in which they flourish. Outside that niche, conditions are less favourable. This leads to negative effects on health, food security and more. This niche is a mean annual temperature below 29 °C. As of May 2023, 60 million people lived outside this niche. With every additional 0.1 degree of warming, 140 million people will be pushed out of it.
Health
Food security
Climate change will affect agriculture and food production around the world. The reasons include the effects of elevated CO2 in the atmosphere. Higher temperatures and altered precipitation and transpiration regimes are also factors. Increased frequency of extreme events and modified weed, pest, and pathogen pressure are other factors. Droughts result in crop failures and the loss of pasture for livestock. Loss and poor growth of livestock cause milk yield and meat production to decrease. The rate of soil erosion is 10–20 times higher than the rate of soil accumulation in agricultural areas that use no-till farming. In areas with tilling it is 100 times higher. Climate change worsens this type of land degradation and desertification.
Climate change is projected to negatively affect all four pillars of food security. It will affect how much food is available. It will also affect how easy food is to access through prices, food quality, and how stable the food system is. Climate change is already affecting the productivity of wheat and other staples.
In many areas, fishery catches are already decreasing because of global warming and changes in biochemical cycles. In combination with overfishing, warming waters decrease the amount of fish in the ocean. Per degree of warming, ocean biomass is expected to decrease by about 5%. Tropical and subtropical oceans are most affected, while there may be more fish in polar waters.
Water security
Water resources can be affected by climate change in various ways. The total amount of freshwater available can change, for instance due to dry spells or droughts. Heavy rainfall and flooding can have an impact on water quality. They can transport pollutants into water bodies through increased surface runoff. In coastal regions, more salt may find its way into water resources due to higher sea levels and more intense storms. Higher temperatures also directly degrade water quality. This is because warm water contains less oxygen. Changes in the water cycle threaten existing and future water infrastructure. It will be harder to plan investments for water infrastructure. This is because there are significant uncertainties about future variability of the water cycle.
Between 1.5 and 2.5 billion people live in areas with regular water security issues. If global warming reaches , water insecurity would affect about twice as many people. Water resources are likely to decrease in most dry subtropical regions and mid-latitudes. But they will increase in high latitudes. However, variable streamflow means even regions with increased water resources can experience additional short-term shortages. In the arid regions of India, China, the US and Africa dry spells and drought are already affecting water availability.
Human settlements
Climate change is particularly likely to affect the Arctic, Africa, small islands, Asian megadeltas and the Middle East regions. Low-latitude, less-developed regions are most at risk of experiencing negative climate change impacts. The ten countries of the Association of Southeast Asian Nations (ASEAN) are among the most vulnerable in the world to the negative effects of climate change. ASEAN's climate mitigation efforts are not in proportion to the climate change threats the region faces.
Impacts from heat
Regions inhabited by a third of the human population could become as hot as the hottest parts of the Sahara within 50 years. This would happen if greenhouse gas emissions continue to grow rapdily without a change in patterns of population growth and without migration. The projected average temperature of above for these regions would be outside the "human temperature niche". This is a range for climate that is biologically suitable for humans. It is based on historical data of mean annual temperatures. The most affected regions have little adaptive capacity.
Increased extreme heat exposure from climate change and the urban heat island effect threatens urban settlements. This is made worse by the loss of shade from urban trees that cannot withstand the heat stress.
In 2019, the Crowther Lab from ETH Zurich paired the climatic conditions of 520 major cities worldwide with the predicted climatic conditions of cities in 2050. It found that 22% of the major cities would have climatic conditions that do not exist in any city today. For instance, 2050 London would have a climate similar to 2019 Melbourne in Australia. Athens and Madrid would be like Fez in Morocco. Nairobi in Kenya would be like Maputo in Mozambique. The Indian city Pune would be like Bamako in Mali and Bamako would be like Niamey in Niger. Brasilia would be like Goiania, both in Brazil.
Low-lying coastal regions
Low-lying cities and other settlements near the sea face multiple simultaneous risks from climate change. They face flooding risks from sea level rise. In addition they may face impacts from more severe storms, ocean acidification, and salt intrusion into the groundwater. Changes like continued development in exposed areas increase the risks that these regions face.
Population density on the coasts is high. Estimates of the number of people at risk of coastal flooding from climate-driven sea level rise vary. Estimates range from 190 million to 300 million. It could even be 640 million in a worst-case scenario related to the instability of the Antarctic ice sheet. People are most affected in the densely-populated low-lying megadeltas of Asia and Africa.
Small island developing states are especially vulnerable. They are likely to experience more intense storm surges, salt water intrusion and coastal destruction. Low-lying small islands in the Pacific, Indian, and Caribbean regions even risk permanent inundation. This would displace their population. On the islands of Fiji, Tonga and western Samoa, migrants from outer islands inhabit low and unsafe areas along the coasts. The entire populations of small atoll nations such as Kiribati, Maldives, the Marshall Islands, and Tuvalu are at risk of being displaced. This could raise issues of statelessness. Several factors increase their vulnerability. These are small size, isolation from other land, low financial resources, and lack of protective infrastructure.
Impacts on societies
Climate change has many impacts on society. It affects health, the availability of drinking water and food, inequality and economic growth. The effects of climate change are often interlinked. They can exacerbate each other as well as existing vulnerabilities. Some areas may become too hot for humans to live in. Climate-related changes or disasters may lead people in some areas to move to other parts of the country or to other countries.
Some scientists describe the effects of climate change, with continuing increases in greenhouse gas emissions, as a "climate emergency" or "climate crisis". Some researchers and activists describe them as an existential threat to civilization. Some define these threats under climate security. The consequences of climate change, and the failure to address it, can distract people from tackling its root causes. This leads to what some researchers have termed a "climate doom loop".
Displacement and migration
Displacement is when people move within a country. Migration is when they move to another country. Some people use the terms interchangeably. Climate change affects displacement in several ways. More frequent and severe weather-related disasters may increase involuntary displacement. These destroy homes and habitats. Climate impacts such as desertification and rising sea levels gradually erode livelihoods. They force communities to abandon traditional homelands. Other forms of migration are adaptive and voluntary. They are based on individual or household decisions. On the other hand, some households may fall into poverty or get poorer due to climate change. This limits their ability to move to less affected areas.
Migration due to climate and weather is usually within countries. But it is long-distance. Slow-onset disasters such as droughts and heat are more likely to cause long-term migration than weather disasters like floods. Migration due to desertification and reduced soil fertility is typically from rural areas in developing countries to towns and cities.
According to the Internal Displacement Monitoring Centre, extreme weather events displaced approximately 30 million people in 2020. Violence and wars displaced approximately 10 million in the same year. There may have been a contribution of climate change to these conflicts. In 2018, the World Bank estimated that climate change will cause internal migration of between 31 and 143 million people by 2050. This would be as they escape crop failures, water scarcity, and sea level rise. The study covered only Sub-Saharan Africa, South Asia, and Latin America.
Conflict
Climate change is unlikely to cause international wars in the foreseeable future. However, climate change can increase the risk for intrastate conflicts, such as civil wars, communal violence, or protests. The IPCC Sixth Assessment Report concludes: "Climate hazards have affected armed conflict within countries (medium confidence), but the influence of climate is small compared to socio-economic, political, and cultural factors (high confidence)."
Climate change can increase conflict risks by causing tensions about scarce resources like food, water and land, by weakening state institutions, by reducing the opportunity costs for impoverished individuals to join armed groups, and by causing tensions related to (climate-induced) migration. Efforts to mitigate or adapt to climate change can also cause conflicts, for instance due to higher food and energy prices or when people are forcibly re-located from vulnerable areas.
Research has shown that climate change is not the most important conflict driver, and that it can only affect conflict risks under certain circumstances. Relevant context factors include agricultural dependence, a history of political instability, poverty, and the political exclusion of ethnic groups. Climate change has thus been described as a "threat multiplier". Yet, an impact of climate change on specific conflicts like the Syrian civil war or the armed conflict in Darfur remains hard to prove.
Social impacts on vulnerable groups
Climate change does not affect people within communities in the same way. It can have a bigger impact on vulnerable groups such as women, the elderly, religious minorities and refugees than on others.
People living with disability. Climate impacts on disabled people have been identified by activists and advocacy groups as well as through the UNHCR adopting a resolution on climate change and the rights of people with disabilities.
People living in poverty: Climate change disproportionally affects poor people in low-income communities and developing countries around the world. Those in poverty have a higher chance of experiencing the ill-effects of climate change, due to their increased exposure and vulnerability. A 2020 World Bank paper estimated that between 32 million to 132 million additional people will be pushed into extreme poverty by 2030 due to climate change.
Women: Climate change increases gender inequality. It reduces women's ability to be financially independent, and has an overall negative impact on the social and political rights of women. This is especially the case in economies that are heavily based on agriculture.
Indigenous peoples: Indigenous communities tend to rely more on the environment for food and other necessities. This makes them more vulnerable to disturbances in ecosystems. Indigenous communities across the globe generally have bigger economic disadvantages than non-indigenous communities. This is due to the oppression they have experienced. These disadvantages include less access to education and jobs and higher rates of poverty. All this makes them more vulnerable to climate change.
Children: The Lancet review on health and climate change lists children among the worst-affected by global warming. Children are 14–44 percent more likely to die from environmental factors.
Possibility of societal collapse
Climate change has long been described as a severe risk to humans. Climate change as an existential threat has emerged as a key theme in the climate movement. People from small island nations also use this theme. There has not been extensive research in this topic. Existential risks are threats that could cause the extinction of humanity or destroy the potential of intelligent life on Earth. Key risks of climate change do not fit that definition. However, some key climate risks do have an impact people's ability to survive. For instance, areas may become too hot to survive, or sea level rise may make it impossible to live at a specific location.
As of October 2024, the possibility of societal collapse became more probable, the number of articles speaking about climate change and societal collapse increased sharply. Leading climate scientists emphasize that "“Climate change is a glaring symptom of a deeper systemic issue: ecological overshoot, [which] is an inherently unstable state that cannot persist indefinitely". To prevent it, they propose phase down fossil fuels, reduce methane emissions, overconsumption, and birth rate, switch to plant-based food, protect and restore ecosystems and adopt an ecological, post-growth economics which includes social justice. Climate change education should be integrated into core curriculums worldwide.
Economic impacts
Economic forecasts of the impact of global warming vary considerably. The impacts are worse if there is insufficient adaptation. Economic modelling may underrate the impact of catastrophic climatic changes. When estimating losses, economists choose a discount rate. This determines how much one prefers to have goods or cash now compared to at a future date. Using a high discount rate may understate economic losses. This is because losses for future generations weigh less heavily.
Economic impacts are bigger the more the temperature rises. Scientists have compared impacts with warming of 1.5 °C (2.7 °F) and a level of 3.66 °C (6.59 °F). They use this higher figure to represent no efforts to stop emissions. They found that total damages at 1.5 °C were 90% less than at 3.66 °C. One study found that global GDP at the end of the century would be 3.5% less if warming is limited to . This study excludes the potential effect of tipping points. Another study found that excluding tipping points underestimates the global economic impact by a factor of two to eight. Another study found that a temperature rise of by 2050 would reduce global GDP by 2.5%–7.5%. By 2100 in this scenario the temperature would rise by . This could reduce global GDP by 30% in the worst case. A 2024 study, which checked the data from the last 120 years, found that climate change has already reduced welfare by 29% and further temperature rise will rise the number to 47%. The temperature rise during the years 1960-2019 alone has cut current GDP per capita by 18%. A 1 degree warming reduces global GDP by 12%. An increase of 3 degrees by 2100, will reduce capital by 50%. The effects are similar to experiencing the 1929 Great Depression permanently. The correct social cost of carbon according to the study is 1065 dollars per tonne of CO2.
Global losses reveal rapidly rising costs due to extreme weather events since the 1970s. Socio-economic factors have contributed to the observed trend of global losses. These factors include population growth and increased wealth. Regional climatic factors also play a role. These include changes in precipitation and flooding events. It is difficult to quantify the relative impact of socio-economic factors and climate change on the observed trend. The trend does suggest social systems are increasing vulnerable to climate change.
Economic inequality
Climate change has contributed to global economic inequality. Wealthy countries in colder regions have felt little overall economic impact from climate change or may have benefited. Poor hotter countries probably grew less than if there had been no global warming.
Highly affected sectors
Climate change has a bigger impact on economic sectors directly affected by weather than on other sectors. It heavily affects agriculture, fisheries and forestry. It also affects the tourism and energy sectors. Agriculture and forestry have suffered economic losses due to droughts and extreme heat. If global warming goes over 1.5 °C, there may be limits to how much tourism and outdoor work can adapt.
In the energy sector, thermal power plants depend on water to cool them. Climate change can increase the likelihood of drought and fresh water shortages. Higher operating temperatures make them less efficient. This reduces their output. Hydropower is affected by changes in the water cycle such as river flows. Diminished river flows can cause power shortages in areas that depend on hydroelectric power. Brazil relies on hydroelectricity. So it is particularly vulnerable. Rising temperatures, lower water flow, and changes in rainfall could reduce total energy production by 7% annually by the end of the century. Climate change affects oil and natural gas infrastructure. This is also vulnerable to the increased risk of disasters such as storms, cyclones, flooding and rising sea levels.
Global warming affects the insurance and financial services sectors. Insurance is an important tool to manage risks. But it is often unavailable to poorer households. Due to climate change, premiums are going up for certain types of insurance, such as flood insurance. Poor adaptation to climate change further widens the gap between what people can afford and the costs of insurance, as risks increase. In 2019 Munich Re said climate change could make home insurance unaffordable for households at or below average incomes.
It is possible that climate change has already begun to affect the shipping sector by impacting the Panama Canal. Lack of rainfall possibly linked to climate change reduced the number of ships passing through the canal per day, from 36 to 22 and by February 2024, it is expected to be 18.
| Physical sciences | Climate change | Earth science |
2119179 | https://en.wikipedia.org/wiki/Climate%20change%20mitigation | Climate change mitigation | Climate change mitigation (or decarbonisation) is action to limit the greenhouse gases in the atmosphere that cause climate change. Climate change mitigation actions include conserving energy and replacing fossil fuels with clean energy sources. Secondary mitigation strategies include changes to land use and removing carbon dioxide (CO2) from the atmosphere. Current climate change mitigation policies are insufficient as they would still result in global warming of about 2.7 °C by 2100, significantly above the 2015 Paris Agreement's goal of limiting global warming to below 2 °C.
Solar energy and wind power can replace fossil fuels at the lowest cost compared to other renewable energy options. The availability of sunshine and wind is variable and can require electrical grid upgrades, such as using long-distance electricity transmission to group a range of power sources. Energy storage can also be used to even out power output, and demand management can limit power use when power generation is low. Cleanly generated electricity can usually replace fossil fuels for powering transportation, heating buildings, and running industrial processes. Certain processes are more difficult to decarbonise, such as air travel and cement production. Carbon capture and storage (CCS) can be an option to reduce net emissions in these circumstances, although fossil fuel power plants with CCS technology is currently a high cost climate change mitigation strategy.
Human land use changes such as agriculture and deforestation cause about 1/4th of climate change. These changes impact how much is absorbed by plant matter and how much organic matter decays or burns to release . These changes are part of the fast carbon cycle, whereas fossil fuels release that was buried underground as part of the slow carbon cycle. Methane is a short lived greenhouse gas that is produced by decaying organic matter and livestock, as well as fossil fuel extraction. Land use changes can also impact precipitation patterns and the reflectivity of the surface of the Earth. It is possible to cut emissions from agriculture by reducing food waste, switching to a more plant-based diet (also referred to as low-carbon diet), and by improving farming processes.
Various policies can encourage climate change mitigation. Carbon pricing systems have been set up that either tax emissions or cap total emissions and trade emission credits. Fossil fuel subsidies can be eliminated in favor of clean energy subsidies, and incentives offered for installing energy efficiency measures or switching to electric power sources. Another issue is overcoming environmental objections when constructing new clean energy sources and making grid modifications. Limiting climate change by reducing greenhouse gas emissions or removing greenhouse gases from the atmosphere could be supplemented by climate technologies such as solar radiation management (or solar geoengineering). Complementary climate change actions, including climate activism, have a focus on political and cultural aspects.
Definitions and scope
Climate change mitigation aims to sustain ecosystems to maintain human civilisation. This requires drastic cuts in greenhouse gas emissions . The Intergovernmental Panel on Climate Change (IPCC) defines mitigation (of climate change) as "a human intervention to reduce emissions or enhance the sinks of greenhouse gases".
It is possible to approach various mitigation measures in parallel. This is because there is no single pathway to limit global warming to 1.5 or 2 °C. There are four types of measures:
Sustainable energy and sustainable transport
Energy conservation, including efficient energy use
Sustainable agriculture and green industrial policy
Enhancing carbon sinks and carbon dioxide removal (CDR), including carbon sequestration
The IPCC defined carbon dioxide removal as "Anthropogenic activities removing carbon dioxide () from the atmosphere and durably storing it in geological, terrestrial, or ocean reservoirs, or in products. It includes existing and potential anthropogenic enhancement of biological or geochemical sinks and direct air carbon dioxide capture and storage (DACCS), but excludes natural uptake not directly caused by human activities."
Emission trends and pledges
Greenhouse gas emissions from human activities strengthen the greenhouse effect. This contributes to climate change. Most is carbon dioxide from burning fossil fuels: coal, oil, and natural gas. Human-caused emissions have increased atmospheric carbon dioxide by about 50% over pre-industrial levels. Emissions in the 2010s averaged a record 56 billion tons (Gt) a year. In 2016, energy for electricity, heat and transport was responsible for 73.2% of GHG emissions. Direct industrial processes accounted for 5.2%, waste for 3.2% and agriculture, forestry and land use for 18.4%.
Electricity generation and transport are major emitters. The largest single source is coal-fired power stations with 20% of greenhouse gas emissions. Deforestation and other changes in land use also emit carbon dioxide and methane. The largest sources of anthropogenic methane emissions are agriculture, and gas venting and fugitive emissions from the fossil-fuel industry. The largest agricultural methane source is livestock. Agricultural soils emit nitrous oxide, partly due to fertilizers. There is now a political solution to the problem of fluorinated gases from refrigerants. This is because many countries have ratified the Kigali Amendment.
Carbon dioxide () is the dominant emitted greenhouse gas. Methane () emissions almost have the same short-term impact. Nitrous oxide (N2O) and fluorinated gases (F-Gases) play a minor role. Livestock and manure produce 5.8% of all greenhouse gas emissions. But this depends on the time frame used to calculate the global warming potential of the respective gas.
Greenhouse gas (GHG) emissions are measured in equivalents. Scientists determine their equivalents from their global warming potential (GWP). This depends on their lifetime in the atmosphere. There are widely used greenhouse gas accounting methods that convert volumes of methane, nitrous oxide and other greenhouse gases to carbon dioxide equivalents. Estimates largely depend on the ability of oceans and land sinks to absorb these gases. Short-lived climate pollutants (SLCPs) persist in the atmosphere for a period ranging from days to 15 years. Carbon dioxide can remain in the atmosphere for millennia. Short-lived climate pollutants include methane, hydrofluorocarbons (HFCs), tropospheric ozone and black carbon.
Scientists increasingly use satellites to locate and measure greenhouse gas emissions and deforestation. Earlier, scientists largely relied on or calculated estimates of greenhouse gas emissions and governments' self-reported data.
Needed emissions cuts
The annual "Emissions Gap Report" by UNEP stated in 2022 that it was necessary to almost halve emissions. "To get on track for limiting global warming to 1.5°C, global annual GHG emissions must be reduced by 45 per cent compared with emissions projections under policies currently in place in just eight years, and they must continue to decline rapidly after 2030, to avoid exhausting the limited remaining atmospheric carbon budget." The report commented that the world should focus on broad-based economy-wide transformations and not incremental change.
In 2022, the Intergovernmental Panel on Climate Change (IPCC) released its Sixth Assessment Report on climate change. It warned that greenhouse gas emissions must peak before 2025 at the latest and decline 43% by 2030 to have a good chance of limiting global warming to 1.5 °C (2.7 °F). Or in the words of Secretary-General of the United Nations António Guterres: "Main emitters must drastically cut emissions starting this year".
Pledges
Climate Action Tracker described the situation on 9 November 2021 as follows. The global temperature will rise by 2.7 °C by the end of the century with current policies and by 2.9 °C with nationally adopted policies. The temperature will rise by 2.4 °C if countries only implement the pledges for 2030. The rise would be 2.1 °C with the achievement of the long-term targets too. Full achievement of all announced targets would mean the rise in global temperature will peak at 1.9 °C and go down to 1.8 °C by the year 2100. Experts gather information about climate pledges in the Global Climate Action Portal - Nazca. The scientific community is checking their fulfilment.
There has not been a definitive or detailed evaluation of most goals set for 2020. But it appears the world failed to meet most or all international goals set for that year.
One update came during the 2021 United Nations Climate Change Conference in Glasgow. The group of researchers running the Climate Action Tracker looked at countries responsible for 85% of greenhouse gas emissions. It found that only four countries or political entities—the EU, UK, Chile and Costa Rica—have published a detailed official policyplan that describes the steps to realise 2030 mitigation targets. These four polities are responsible for 6% of global greenhouse gas emissions.
In 2021 the US and EU launched the Global Methane Pledge to cut methane emissions by 30% by 2030. The UK, Argentina, Indonesia, Italy and Mexico joined the initiative. Ghana and Iraq signaled interest in joining. A White House summary of the meeting noted those countries represent six of the top 15 methane emitters globally. Israel also joined the initiative.
Low-carbon energy
The energy system includes the delivery and use of energy. It is the main emitter of carbon dioxide (). Rapid and deep reductions in the carbon dioxide and other greenhouse gas emissions from the energy sector are necessary to limit global warming to well below 2 °C. IPCC recommendations include reducing fossil fuel consumption, increasing production from low- and zero carbon energy sources, and increasing use of electricity and alternative energy carriers.
Nearly all scenarios and strategies involve a major increase in the use of renewable energy in combination with increased energy efficiency measures. It will be necessary to accelerate the deployment of renewable energy six-fold from 0.25% annual growth in 2015 to 1.5% to keep global warming under 2 °C.
The competitiveness of renewable energy is a key to a rapid deployment. In 2020, onshore wind and solar photovoltaics were the cheapest source for new bulk electricity generation in many regions. Renewables may have higher storage costs but non-renewables may have higher clean-up costs. A carbon price can increase the competitiveness of renewable energy.
Solar and wind energy
Wind and sun can provide large amounts of low-carbon energy at competitive production costs. The IPCC estimates that these two mitigation options have the largest potential to reduce emissions before 2030 at low cost.
Solar photovoltaics (PV) has become the cheapest way to generate electricity in many regions of the world. The growth of photovoltaics has been close to exponential. It has about doubled every three years since the 1990s. A different technology is concentrated solar power (CSP). This uses mirrors or lenses to concentrate a large area of sunlight on to a receiver. With CSP, the energy can be stored for a few hours. This provides supply in the evening. Solar water heating doubled between 2010 and 2019.
Regions in the higher northern and southern latitudes have the greatest potential for wind power. Offshore wind farms are more expensive. But offshore units deliver more energy per installed capacity with less fluctuations. In most regions, wind power generation is higher in the winter when PV output is low. For this reason, combinations of wind and solar power lead to better-balanced systems.
Other renewables
Other well-established renewable energy forms include hydropower, bioenergy and geothermal energy.
Hydroelectricity is electricity generated by hydropower and plays a leading role in countries like Brazil, Norway and China. but there are geographical limits and environmental issues. Tidal power can be used in coastal regions.
Bioenergy can provide energy for electricity, heat and transport. Bioenergy, in particular biogas, can provide dispatchable electricity generation. While burning plant-derived biomass releases , the plants withdraw from the atmosphere while they grow. The technologies for producing, transporting and processing a fuel have a significant impact on the lifecycle emissions of the fuel. For example, aviation is starting to use renewable biofuels.
Geothermal power is electrical power generated from geothermal energy. Geothermal electricity generation is currently used in 26 countries. Geothermal heating is in use in 70 countries.
Integrating variable renewable energy
Wind and solar power production does not consistently match demand. To deliver reliable electricity from variable renewable energy sources such as wind and solar, electrical power systems must be flexible. Most electrical grids were constructed for non-intermittent energy sources such as coal-fired power plants. The integration of larger amounts of solar and wind energy into the grid requires a change of the energy system; this is necessary to ensure that the supply of electricity matches demand.
There are various ways to make the electricity system more flexible. In many places, wind and solar generation are complementary on a daily and a seasonal scale. There is more wind during the night and in winter when solar energy production is low. Linking different geographical regions through long-distance transmission lines also makes it possible to reduce variability. It is possible to shift energy demand in time. Energy demand management and the use of smart grids make it possible to match the times when variable energy production is highest. Sector coupling can provide further flexibility. This involves coupling the electricity sector to the heat and mobility sector via power-to-heat-systems and electric vehicles.
Energy storage helps overcome barriers to intermittent renewable energy. The most commonly used and available storage method is pumped-storage hydroelectricity. This requires locations with large differences in height and access to water. Batteries are also in wide use. They typically store electricity for short periods. Batteries have low energy density. This and their cost makes them impractical for the large energy storage necessary to balance inter-seasonal variations in energy production. Some locations have implemented pumped hydro storage with capacity for multi-month usage.
Nuclear power
Nuclear power could complement renewables for electricity. On the other hand, environmental and security risks could outweigh the benefits.
The construction of new nuclear reactors currently takes about 10 years. This is much longer than scaling up the deployment of wind and solar. And this timing gives rise to credit risks. However nuclear may be much cheaper in China. China is building a significant number of new power plants. the cost of extending nuclear power plant lifetimes is competitive with other electricity generation technologies if long term costs for nuclear waste disposal are excluded from the calculation. There is also no sufficient financial insurance for nuclear accidents.
Replacing coal with natural gas
Demand reduction
Reducing demand for products and services that cause greenhouse gas emissions can help in mitigating climate change. One is to reduce demand by behavioural and cultural changes, for example by making changes in diet, especially the decision to reduce meat consumption, an effective action individuals take to fight climate change. Another is by reducing the demand by improving infrastructure, by building a good public transport network, for example. Lastly, changes in end-use technology can reduce energy demand. For instance a well-insulated house emits less than a poorly-insulated house.
Mitigation options that reduce demand for products or services help people make personal choices to reduce their carbon footprint. This could be in their choice of transport or food. So these mitigation options have many social aspects that focus on demand reduction; they are therefore demand-side mitigation actions. For example, people with high socio-economic status often cause more greenhouse gas emissions than those from a lower status. If they reduce their emissions and promote green policies, these people could become low-carbon lifestyle role models. However, there are many psychological variables that influence consumers. These include awareness and perceived risk.
Government policies can support or hinder demand-side mitigation options. For example, public policy can promote circular economy concepts which would support climate change mitigation. Reducing greenhouse gas emissions is linked to the sharing economy.
There is a debate regarding the correlation of economic growth and emissions. It seems economic growth no longer necessarily means higher emissions.
Energy conservation and efficiency
Global primary energy demand exceeded 161,000 terawatt hours (TWh) in 2018. This refers to electricity, transport and heating including all losses. In transport and electricity production, fossil fuel usage has a low efficiency of less than 50%. Large amounts of heat in power plants and in motors of vehicles go to waste. The actual amount of energy consumed is significantly lower at 116,000 TWh.
Energy conservation is the effort made to reduce the consumption of energy by using less of an energy service. One way is to use energy more efficiently. This means using less energy than before to produce the same service. Another way is to reduce the amount of service used. An example of this would be to drive less. Energy conservation is at the top of the sustainable energy hierarchy. When consumers reduce wastage and losses they can conserve energy. The upgrading of technology as well as the improvements to operations and maintenance can result in overall efficiency improvements.
Efficient energy use (or energy efficiency) is the process of reducing the amount of energy required to provide products and services. Improved energy efficiency in buildings ("green buildings"), industrial processes and transportation could reduce the world's energy needs in 2050 by one third. This would help reduce global emissions of greenhouse gases. For example, insulating a building allows it to use less heating and cooling energy to achieve and maintain thermal comfort. Improvements in energy efficiency are generally achieved by adopting a more efficient technology or production process. Another way is to use commonly accepted methods to reduce energy losses.
Lifestyle changes
Individual action on climate change can include personal choices in many areas. These include diet, travel, household energy use, consumption of goods and services, and family size. People who wish to reduce their carbon footprint can take high-impact actions such as avoiding frequent flying and petrol-fuelled cars, eating mainly a plant-based diet, having fewer children, using clothes and electrical products for longer, and electrifying homes. These approaches are more practical for people in high-income countries with high-consumption lifestyles. Naturally, it is more difficult for those with lower income statuses to make these changes. This is because choices like electric-powered cars may not be available. Excessive consumption is more to blame for climate change than population increase. High-consumption lifestyles have a greater environmental impact, with the richest 10% of people emitting about half the total lifestyle emissions.
Dietary change
Some scientists say that avoiding meat and dairy foods is the single biggest way an individual can reduce their environmental impact. The widespread adoption of a vegetarian diet could cut food-related greenhouse gas emissions by 63% by 2050. China introduced new dietary guidelines in 2016 which aim to cut meat consumption by 50% and thereby reduce greenhouse gas emissions by 1Gt per year by 2030. Overall, food accounts for the largest share of consumption-based greenhouse gas emissions. It is responsible for nearly 20% of the global carbon footprint. Almost 15% of all anthropogenic greenhouse gas emissions have been attributed to the livestock sector.
A shift towards plant-based diets would help to mitigate climate change. In particular, reducing meat consumption would help to reduce methane emissions. If high-income nations switched to a plant-based diet, vast amounts of land used for animal agriculture could be allowed to return to their natural state. This in turn has the potential to sequester 100 billion tonnes of by the end of the century. A comprehensive analysis found that plant based diets reduce emissions, water pollution and land use significantly (by 75%), while reducing the destruction of wildlife and usage of water.
Family size
Population growth has resulted in higher greenhouse gas emissions in most regions, particularly Africa. However, economic growth has a bigger effect than population growth. Rising incomes, changes in consumption and dietary patterns, as well as population growth, cause pressure on land and other natural resources. This leads to more greenhouse gas emissions and fewer carbon sinks. Some scholars have argued that humane policies to slow population growth should be part of a broad climate response together with policies that end fossil fuel use and encourage sustainable consumption. Advances in female education and reproductive health, especially voluntary family planning, can contribute to reducing population growth.
Preserving and enhancing carbon sinks
An important mitigation measure is "preserving and enhancing carbon sinks". This refers to the management of Earth's natural carbon sinks in a way that preserves or increases their capability to remove CO2 from the atmosphere and to store it durably. Scientists call this process also carbon sequestration. In the context of climate change mitigation, the IPCC defines a sink as "Any process, activity or mechanism which removes a greenhouse gas, an aerosol or a precursor of a greenhouse gas from the atmosphere". Globally, the two most important carbon sinks are vegetation and the ocean.
To enhance the ability of ecosystems to sequester carbon, changes are necessary in agriculture and forestry. Examples are preventing deforestation and restoring natural ecosystems by reforestation. Scenarios that limit global warming to 1.5 °C typically project the large-scale use of carbon dioxide removal methods over the 21st century. There are concerns about over-reliance on these technologies, and their environmental impacts. But ecosystem restoration and reduced conversion are among the mitigation tools that can yield the most emissions reductions before 2030.
Land-based mitigation options are referred to as "AFOLU mitigation options" in the 2022 IPCC report on mitigation. The abbreviation stands for "agriculture, forestry and other land use" The report described the economic mitigation potential from relevant activities around forests and ecosystems as follows: "the conservation, improved management, and restoration of forests and other ecosystems (coastal wetlands, peatlands, savannas and grasslands)". A high mitigation potential is found for reducing deforestation in tropical regions. The economic potential of these activities has been estimated to be 4.2 to 7.4 gigatonnes of carbon dioxide equivalent (GtCO2 -eq) per year.
Forests
Conservation
The Stern Review on the economics of climate change stated in 2007 that curbing deforestation was a highly cost-effective way of reducing greenhouse gas emissions. About 95% of deforestation occurs in the tropics, where clearing of land for agriculture is one of the main causes. One forest conservation strategy is to transfer rights over land from public ownership to its indigenous inhabitants. Land concessions often go to powerful extractive companies. Conservation strategies that exclude and even evict humans, called fortress conservation, often lead to more exploitation of the land. This is because the native inhabitants turn to work for extractive companies to survive.
Proforestation is promoting forests to capture their full ecological potential. This is a mitigation strategy as secondary forests that have regrown in abandoned farmland are found to have less biodiversity than the original old-growth forests. Original forests store 60% more carbon than these new forests. Strategies include rewilding and establishing wildlife corridors.
Afforestation and reforestation
Afforestation is the establishment of trees where there was previously no tree cover. Scenarios for new plantations covering up to 4000 million hectares (Mha) (6300 x 6300 km) suggest cumulative carbon storage of more than 900 GtC (2300 Gt) until 2100. But they are not a viable alternative to aggressive emissions reduction. This is because the plantations would need to be so large they would eliminate most natural ecosystems or reduce food production. One example is the Trillion Tree Campaign. However, preserving biodiversity is also important and for example not all grasslands are suitable for conversion into forests. Grasslands can even turn from carbon sinks to carbon sources.
Reforestation is the restocking of existing depleted forests or in places where there were recently forests. Reforestation could save at least 1GtCO2 per year, at an estimated cost of $5–15 per tonne of carbon dioxide (tCO2). Restoring all degraded forests all over the world could capture about 205 GtC (750 Gt). With increased intensive agriculture and urbanization, there is an increase in the amount of abandoned farmland. By some estimates, for every acre of original old-growth forest cut down, more than 50 acres of new secondary forests are growing. In some countries, promoting regrowth on abandoned farmland could offset years of emissions.
Planting new trees can be expensive and a risky investment. For example, about 80 percent of planted trees in the Sahel die within two years. Reforestation has higher carbon storage potential than afforestation. Even long-deforested areas still contain an "underground forest" of living roots and tree stumps. Helping native species sprout naturally is cheaper than planting new trees and they are more likely to survive. This could include pruning and coppicing to accelerate growth. This also provides woodfuel, which is otherwise a major source of deforestation. Such practices, called farmer-managed natural regeneration, are centuries old but the biggest obstacle towards implementation is ownership of the trees by the state. The state often sells timber rights to businesses which leads to locals uprooting seedlings because they see them as a liability. Legal aid for locals and changes to property law such as in Mali and Niger have led to significant changes. Scientists describe them as the largest positive environmental transformation in Africa. It is possible to discern from space the border between Niger and the more barren land in Nigeria, where the law has not changed.
Soils
There are many measures to increase soil carbon. This makes it complex and hard to measure and account for. One advantage is that there are fewer trade-offs for these measures than for BECCS or afforestation, for example.
Globally, protecting healthy soils and restoring the soil carbon sponge could remove 7.6 billion tonnes of carbon dioxide from the atmosphere annually. This is more than the annual emissions of the US. Trees capture while growing above ground and exuding larger amounts of carbon below ground. Trees contribute to the building of a soil carbon sponge. Carbon formed above ground is released as immediately when wood is burned. If dead wood remains untouched, only some of the carbon returns to the atmosphere as decomposition proceeds.
Farming can deplete soil carbon and render soil incapable of supporting life. However, conservation farming can protect carbon in soils, and repair damage over time. The farming practice of cover crops is a form of carbon farming. Methods that enhance carbon sequestration in soil include no-till farming, residue mulching and crop rotation. Scientists have described the best management practices for European soils to increase soil organic carbon. These are conversion of arable land to grassland, straw incorporation, reduced tillage, straw incorporation combined with reduced tillage, ley cropping system and cover crops.
Another mitigation option is the production of biochar and its storage in soils This is the solid material that remains after the pyrolysis of biomass. Biochar production releases half of the carbon from the biomass—either released into the atmosphere or captured with CCS—and retains the other half in the stable biochar. It can endure in soil for thousands of years. Biochar may increase the soil fertility of acidic soils and increase agricultural productivity. During production of biochar, heat is released which may be used as bioenergy.
Wetlands
Wetland restoration is an important mitigation measure. It has moderate to great mitigation potential on a limited land area with low trade-offs and costs. Wetlands perform two important functions in relation to climate change. They can sequester carbon, converting carbon dioxide to solid plant material through photosynthesis. They also store and regulate water. Wetlands store about 45 million tonnes of carbon per year globally.
Some wetlands are a significant source of methane emissions. Some also emit nitrous oxide. Peatland globally covers just 3% of the land's surface. But it stores up to 550 gigatonnes (Gt) of carbon. This represents 42% of all soil carbon and exceeds the carbon stored in all other vegetation types, including the world's forests. The threat to peatlands includes draining the areas for agriculture. Another threat is cutting down trees for lumber, as the trees help hold and fix the peatland. Additionally, peat is often sold for compost. It is possible to restore degraded peatlands by blocking drainage channels in the peatland, and allowing natural vegetation to recover.
Mangroves, salt marshes and seagrasses make up the majority of the ocean's vegetated habitats. They only equal 0.05% of the plant biomass on land. But they store carbon 40 times faster than tropical forests. Bottom trawling, dredging for coastal development and fertilizer runoff have damaged coastal habitats. Notably, 85% of oyster reefs globally have been removed in the last two centuries. Oyster reefs clean the water and help other species thrive. This increases biomass in that area. In addition, oyster reefs mitigate the effects of climate change by reducing the force of waves from hurricanes. They also reduce the erosion from rising sea levels. Restoration of coastal wetlands is thought to be more cost-effective than restoration of inland wetlands.
Deep ocean
These options focus on the carbon which ocean reservoirs can store. They include ocean fertilization, ocean alkalinity enhancement or enhanced weathering. The IPCC found in 2022 ocean-based mitigation options currently have only limited deployment potential. But it assessed that their future mitigation potential is large. It found that in total, ocean-based methods could remove 1–100 Gt of per year. Their costs are in the order of US$40–500 per tonne of . Most of these options could also help to reduce ocean acidification. This is the drop in pH value caused by increased atmospheric CO2 concentrations.
Blue carbon management is another type of ocean-based biological carbon dioxide removal (CDR). It can involve land-based as well as ocean-based measures. The term usually refers to the role that tidal marshes, mangroves and seagrasses can play in carbon sequestration. Some of these efforts can also take place in deep ocean waters. This is where the vast majority of ocean carbon is held. These ecosystems can contribute to climate change mitigation and also to ecosystem-based adaptation. Conversely, when blue carbon ecosystems are degraded or lost they release carbon back to the atmosphere. There is increasing interest in developing blue carbon potential. Scientists have found that in some cases these types of ecosystems remove far more carbon per area than terrestrial forests. However, the long-term effectiveness of blue carbon as a carbon dioxide removal solution remains under discussion.
Enhanced weathering
Enhanced weathering could remove 2–4 Gt of per year. This process aims to accelerate natural weathering by spreading finely ground silicate rock, such as basalt, onto surfaces. This speeds up chemical reactions between rocks, water, and air. It removes carbon dioxide from the atmosphere, permanently storing it in solid carbonate minerals or ocean alkalinity. Cost estimates are in the US$50–200 per tonne range of .
Other methods to capture and store CO2
In addition to traditional land-based methods to remove carbon dioxide (CO2) from the air, other technologies are under development. These could reduce CO2 emissions and lower existing atmospheric CO2 levels. Carbon capture and storage (CCS) is a method to mitigate climate change by capturing CO2 from large point sources, such as cement factories or biomass power plants. It then stores it away safely instead of releasing it into the atmosphere. The IPCC estimates that the costs of halting global warming would double without CCS.
Bioenergy with carbon capture and storage (BECCS) expands on the potential of CCS and aims to lower atmospheric CO2 levels. This process uses biomass grown for bioenergy. The biomass yields energy in useful forms such as electricity, heat, biofuels, etc. through consumption of the biomass via combustion, fermentation, or pyrolysis. The process captures the CO2 that was extracted from the atmosphere when it grew. It then stores it underground or via land application as biochar. This effectively removes it from the atmosphere. This makes BECCS a negative emissions technology (NET).
Scientists estimated the potential range of negative emissions from BECCS in 2018 as 0–22 Gt per year. , BECCS was capturing approximately 2 million tonnes per year of CO2 annually. The cost and availability of biomass limits wide deployment of BECCS. BECCS currently forms a big part of achieving climate targets beyond 2050 in modelling, such as by the Integrated Assessment Models (IAMs) associated with the IPCC process. But many scientists are sceptical due to the risk of loss of biodiversity.
Direct air capture is a process of capturing directly from the ambient air. This is in contrast to CCS which captures carbon from point sources. It generates a concentrated stream of for sequestration, utilization or production of carbon-neutral fuel and windgas. Artificial processes vary, and there are concerns about the long-term effects of some of these processes.
Mitigation by sector
Buildings
The building sector accounts for 23% of global energy-related emissions. About half of the energy is used for space and water heating. Building insulation can reduce the primary energy demand significantly. Heat pump loads may also provide a flexible resource that can participate in demand response to integrate variable renewable resources into the grid. Solar water heating uses thermal energy directly. Sufficiency measures include moving to smaller houses when the needs of households change, mixed use of spaces and the collective use of devices. Planners and civil engineers can construct new buildings using passive solar building design, low-energy building, or zero-energy building techniques. In addition, it is possible to design buildings that are more energy-efficient to cool by using lighter-coloured, more reflective materials in the development of urban areas.
Heat pumps efficiently heat buildings, and cool them by air conditioning. A modern heat pump typically transports around three to five times more thermal energy than electrical energy consumed. The amount depends on the coefficient of performance and the outside temperature.
Refrigeration and air conditioning account for about 10% of global emissions caused by fossil fuel-based energy production and the use of fluorinated gases. Alternative cooling systems, such as passive cooling building design and passive daytime radiative cooling surfaces, can reduce air conditioning use. Suburbs and cities in hot and arid climates can significantly reduce energy consumption from cooling with daytime radiative cooling.
Energy consumption for cooling is likely to rise significantly due to increasing heat and availability of devices in poorer countries. Of the 2.8 billion people living in the hottest parts of the world, only 8% currently have air conditioners, compared with 90% of people in the US and Japan. Adoption of air conditioners typically increases in warmer areas at above $10,000 annual household income. By combining energy efficiency improvements and decarbonising electricity for air conditioning with the transition away from super-polluting refrigerants, the world could avoid cumulative greenhouse gas emissions of up to 210–460 Gt-eq over the next four decades. A shift to renewable energy in the cooling sector comes with two advantages: Solar energy production with mid-day peaks corresponds with the load required for cooling and additionally, cooling has a large potential for load management in the electric grid.
Urban planning
Cities emitted 28 GtCO2-eq in 2020 of combined CO2 and emissions. This was from producing and consuming goods and services. Climate-smart urban planning aims to reduce sprawl to reduce the distance travelled. This lowers emissions from transportation. Switching from cars by improving walkability and cycling infrastructure is beneficial to a country's economy as a whole.
Urban forestry, lakes and other blue and green infrastructure can reduce emissions directly and indirectly by reducing energy demand for cooling. Methane emissions from municipal solid waste can be reduced by segregation, composting, and recycling.
Transport
Transportation accounts for 15% of emissions worldwide. Increasing the use of public transport, low-carbon freight transport and cycling are important components of transport decarbonisation.
Electric vehicles and environmentally friendly rail help to reduce the consumption of fossil fuels. In most cases, electric trains are more efficient than air transport and truck transport. Other efficiency means include improved public transport, smart mobility, carsharing and electric hybrids. Fossil-fuel for passenger cars can be included in emissions trading. Furthermore, moving away from a car-dominated transport system towards low-carbon advanced public transport system is important.
Heavyweight, large personal vehicles (such as cars) require a lot of energy to move and take up much urban space. Several alternatives modes of transport are available to replace these. The European Union has made smart mobility part of its European Green Deal. In smart cities, smart mobility is also important.
The World Bank is helping lower income countries buy electric buses. Their purchase price is higher than diesel buses. But lower running costs and health improvements due to cleaner air can offset this higher price.
Between one quarter and three quarters of cars on the road by 2050 are forecast to be electric vehicles. Hydrogen may be a solution for long-distance heavy freight trucks, if batteries alone are too heavy.
Shipping
In the shipping industry, the use of liquefied natural gas (LNG) as a marine bunker fuel is driven by emissions regulations. Ship operators must switch from heavy fuel oil to more expensive oil-based fuels, implement costly flue gas treatment technologies or switch to LNG engines. Methane slip, when gas leaks unburned through the engine, lowers the advantages of LNG. Maersk, the world's biggest container shipping line and vessel operator, warns of stranded assets when investing in transitional fuels like LNG. The company lists green ammonia as one of the preferred fuel types of the future. It has announced the first carbon-neutral vessel on the water by 2023, running on carbon-neutral methanol. Cruise operators are trialling partially hydrogen-powered ships.
Hybrid and all electric ferries are suitable for short distances. Norway's goal is an all electric fleet by 2025.
Air transport
Jet airliners contribute to climate change by emitting carbon dioxide, nitrogen oxides, contrails and particulates. Their radiative forcing is estimated at 1.3–1.4 that of alone, excluding induced cirrus cloud. In 2018, global commercial operations generated 2.4% of all emissions.
The aviation industry has become more fuel efficient. But overall emissions have risen as the volume of air travel has increased. By 2020, aviation emissions were 70% higher than in 2005 and they could grow by 300% by 2050.
It is possible to reduce aviation's environmental footprint by better fuel economy in aircraft. Optimising flight routes to lower non- effects on climate from nitrogen oxides, particulates or contrails can also help. Aviation biofuel, carbon emission trading and carbon offsetting, part of the 191 nation ICAO's Carbon Offsetting and Reduction Scheme for International Aviation (CORSIA), can lower emissions. Short-haul flight bans, train connections, personal choices and taxation on flights can lead to fewer flights. Hybrid electric aircraft and electric aircraft or hydrogen-powered aircraft may replace fossil fuel-powered aircraft.
Experts expect emissions from aviation to rise in most projections, at least until 2040. They currently amount to 180 Mt of or 11% of transport emissions. Aviation biofuel and hydrogen can only cover a small proportion of flights in the coming years. Experts expect hybrid-driven aircraft to start commercial regional scheduled flights after 2030. Battery-powered aircraft are likely to enter the market after 2035. Under CORSIA, flight operators can purchase carbon offsets to cover their emissions above 2019 levels. CORSIA will be compulsory from 2027.
Agriculture, forestry and land use
Almost 20% of greenhouse gas emissions come from the agriculture and forestry sector. To significantly reduce these emissions, annual investments in the agriculture sector need to increase to $260 billion by 2030. The potential benefits from these investments are estimated at about $4.3 trillion by 2030, offering a substantial economic return of 16-to-1.
Mitigation measures in the food system can be divided into four categories. These are demand-side changes, ecosystem protections, mitigation on farms, and mitigation in supply chains. On the demand side, limiting food waste is an effective way to reduce food emissions. Changes to a diet less reliant on animal products such as plant-based diets are also effective.
With 21% of global methane emissions, cattle are a major driver of global warming. When rainforests are cut and the land is converted for grazing, the impact is even higher. In Brazil, producing 1 kg of beef can result in the emission of up to 335 kg CO2-eq. Other livestock, manure management and rice cultivation also emit greenhouse gases, in addition to fossil fuel combustion in agriculture.
Important mitigation options for reducing the greenhouse gas emissions from livestock include genetic selection, introduction of methanotrophic bacteria into the rumen, vaccines, feeds, diet modification and grazing management. Other options are diet changes towards ruminant-free alternatives, such as milk substitutes and meat analogues. Non-ruminant livestock, such as poultry, emit far fewer GHGs.
It is possible to cut methane emissions in rice cultivation by improved water management, combining dry seeding and one drawdown, or executing a sequence of wetting and drying. This results in emission reductions of up to 90% compared to full flooding and even increased yields.
Industry
Industry is the largest emitter of greenhouse gases when direct and indirect emissions are included. Electrification can reduce emissions from industry. Green hydrogen can play a major role in energy-intensive industries for which electricity is not an option. Further mitigation options involve the steel and cement industry, which can switch to a less polluting production process. Products can be made with less material to reduce emission-intensity and industrial processes can be made more efficient. Finally, circular economy measures reduce the need for new materials. This also saves on emissions that would have been released from the mining of collecting of those materials.
The decarbonisation of cement production requires new technologies, and therefore investment in innovation. Bioconcrete is one possibility to reduce emissions. But no technology for mitigation is yet mature. So CCS will be necessary at least in the short-term.
Another sector with a significant carbon footprint is the steel sector, which is responsible for about 7% of global emissions. Emissions can be reduced by using electric arc furnaces to melt and recycle scrap steel. To produce virgin steel without emissions, blast furnaces could be replaced by hydrogen direct reduced iron and electric arc furnaces. Alternatively, carbon capture and storage solutions can be used.
Coal, gas and oil production often come with significant methane leakage. In the early 2020s some governments recognized the scale of the problem and introduced regulations. Methane leaks at oil and gas wells and processing plants are cost-effective to fix in countries which can easily trade gas internationally. There are leaks in countries where gas is cheap; such as Iran, Russia, and Turkmenistan. Nearly all this can be stopped by replacing old components and preventing routine flaring. Coalbed methane may continue leaking even after the mine has been closed. But it can be captured by drainage and/or ventilation systems. Fossil fuel firms do not always have financial incentives to tackle methane leakage.
Co-benefits
Co-benefits of climate change mitigation, also often referred to as ancillary benefits, were firstly dominated in the scientific literature by studies that describe how lower GHG emissions lead to better air quality and consequently impact human health positively. The scope of co-benefits research expanded to its economic, social, ecological and political implications.
Positive secondary effects that occur from climate mitigation and adaptation measures have been mentioned in research since the 1990s. The IPCC first mentioned the role of co-benefits in 2001, followed by its fourth and fifth assessment cycle stressing improved working environment, reduced waste, health benefits and reduced capital expenditures. In the early 2000s the OECD was further fostering its efforts in promoting ancillary benefits.
The IPCC pointed out in 2007: "Co-benefits of GHG mitigation can be an important decision criteria in analyses carried out by policy-makers, but they are often neglected" and added that the co-benefits are "not quantified, monetised or even identified by businesses and decision-makers". Appropriate consideration of co-benefits can greatly "influence policy decisions concerning the timing and level of mitigation action", and there can be "significant advantages to the national economy and technical innovation".
An analysis of climate action in the UK found that public health benefits are a major component of the total benefits derived from climate action.
Employment and economic development
Co-benefits can positively impact employment, industrial development, states' energy independence and energy self-consumption. The deployment of renewable energies can foster job opportunities. Depending on the country and deployment scenario, replacing coal power plants with renewable energy can more than double the number of jobs per average MW capacity. Investments in renewable energies, especially in solar- and wind energy, can boost the value of production. Countries which rely on energy imports can enhance their energy independence and ensure supply security by deploying renewables. National energy generation from renewables lowers the demand for fossil fuel imports which scales up annual economic saving.
The European Commission forecasts a shortage of 180,000 skilled workers in hydrogen production and 66,000 in solar photovoltaic power by 2030.
Energy security
A higher share of renewables can additionally lead to more energy security. Socioeconomic co-benefits have been analysed such as energy access in rural areas and improved rural livelihoods. Rural areas which are not fully electrified can benefit from the deployment of renewable energies. Solar-powered mini-grids can remain economically viable, cost-competitive and reduce the number of power cuts. Energy reliability has additional social implications: stable electricity improves the quality of education.
The International Energy Agency (IEA) spelled out the "multiple benefits approach" of energy efficiency while the International Renewable Energy Agency (IRENA) operationalised the list of co-benefits of the renewable energy sector.
Health and well-being
The health benefits from climate change mitigation are significant. Potential measures can not only mitigate future health impacts from climate change but also improve health directly. Climate change mitigation is interconnected with various health co-benefits, such as those from reduced air pollution. Air pollution generated by fossil fuel combustion is both a major driver of global warming and the cause of a large number of annual deaths. Some estimates are as high as excess deaths during 2018. A 2023 study estimated that fossil fuels kill over 5 million people each year, as of 2019, by causing diseases such as heart attack, stroke and chronic obstructive pulmonary disease. Particulate air pollution kills by far the most, followed by ground-level ozone.
Mitigation policies can also promote healthier diets such as less red meat, more active lifestyles, and increased exposure to green urban spaces. Access to urban green spaces provides benefits to mental health as well. The increased use of green and blue infrastructure can reduce the urban heat island effect. This reduces heat stress on people.
Climate change adaptation
Some mitigation measures have co-benefits in the area of climate change adaptation. This is for example the case for many nature-based solutions. Examples in the urban context include urban green and blue infrastructure which provide mitigation as well as adaptation benefits. This can be in the form of urban forests and street trees, green roofs and walls, urban agriculture and so forth. The mitigation is achieved through the conservation and expansion of carbon sinks and reduced energy use of buildings. Adaptation benefits come for example through reduced heat stress and flooding risk.
Negative side effects
Mitigation measures can also have negative side effects and risks. In agriculture and forestry, mitigation measures can affect biodiversity and ecosystem functioning. In renewable energy, mining for metals and minerals can increase threats to conservation areas. There is some research into ways to recycle solar panels and electronic waste. This would create a source for materials so there is no need to mine them.
Scholars have found that discussions about risks and negative side effects of mitigation measures can lead to deadlock or the feeling that there are insuperable barriers to taking action.
Costs and funding
Several factors affect mitigation cost estimates. One is the baseline. This is a reference scenario that the alternative mitigation scenario is compared with. Others are the way costs are modelled, and assumptions about future government policy. Cost estimates for mitigation for specific regions depend on the quantity of emissions allowed for that region in future, as well as the timing of interventions.
Mitigation costs will vary according to how and when emissions are cut. Early, well-planned action will minimize the costs. Globally, the benefits of keeping warming under 2 °C exceed the costs, which according to The Economist are affordable.
Economists estimate the cost of climate change mitigation at between 1% and 2% of GDP. While this is a large sum, it is still far less than the subsidies governments provide to the ailing fossil fuel industry. The International Monetary Fund estimated this at more than $5 trillion per year.
Another estimate says that financial flows for climate mitigation and adaptation are going to be over $800 billion per year. These financial requirements are predicted to exceed $4 trillion per year by 2030.
Globally, limiting warming to 2 °C may result in higher economic benefits than economic costs. The economic repercussions of mitigation vary widely across regions and households, depending on policy design and level of international cooperation. Delayed global cooperation increases policy costs across regions, especially in those that are relatively carbon intensive at present. Pathways with uniform carbon values show higher mitigation costs in more carbon-intensive regions, in fossil-fuels exporting regions and in poorer regions. Aggregate quantifications expressed in GDP or monetary terms undervalue the economic effects on households in poorer countries. The actual effects on welfare and well-being are comparatively larger.
Cost–benefit analysis may be unsuitable for analysing climate change mitigation as a whole. But it is still useful for analysing the difference between a 1.5 °C target and 2 °C. One way of estimating the cost of reducing emissions is by considering the likely costs of potential technological and output changes. Policymakers can compare the marginal abatement costs of different methods to assess the cost and amount of possible abatement over time. The marginal abatement costs of the various measures will differ by country, by sector, and over time.
Eco-tariffs on only imports contribute to reduced global export competitiveness and to deindustrialization.
Avoided costs of climate change effects
It is possible to avoid some of the costs of the effects of climate change by limiting climate change. According to the Stern Review, inaction can be as high as the equivalent of losing at least 5% of global gross domestic product (GDP) each year, now and forever. This can be up to 20% of GDP or more when including a wider range of risks and impacts. But mitigating climate change will only cost about 2% of GDP. Also it may not be a good idea from a financial perspective to delay significant reductions in greenhouse gas emissions.
Mitigation solutions are often evaluated in terms of costs and greenhouse gas reduction potentials. This fails to take into account the direct effects on human well-being.
Distributing emissions abatement costs
Mitigation at the speed and scale required to limit warming to 2 °C or below implies deep economic and structural changes. These raise multiple types of distributional concerns across regions, income classes and sectors.
There have been different proposals on how to allocate responsibility for cutting emissions. These include egalitarianism, basic needs according to a minimum level of consumption, proportionality and the polluter-pays principle. A specific proposal is "equal per capita entitlements". This approach has two categories. In the first category, emissions are allocated according to national population. In the second category, emissions are allocated in a way that attempts to account for historical or cumulative emissions.
Funding
In order to reconcile economic development with mitigating carbon emissions, developing countries need particular support. This would be both financial and technical. The IPCC found that accelerated support would also tackle inequities in financial and economic vulnerability to climate change. One way to achieve this is the Kyoto Protocol's Clean Development Mechanism (CDM).
Policies
National policies
Climate change mitigation policies can have a large and complex impact on the socio-economic status of individuals and countries This can be both positive and negative. It is important to design policies well and make them inclusive. Otherwise climate change mitigation measures can impose higher financial costs on poor households.
An evaluation was conducted on 1,500 climate policy interventions made between 1998 and 2022. The interventions took place in 41 countries and across 6 continents, which together contributed 81% of the world's total emissions as of 2019. The evaluation found 63 successful interventions that resulted in significant emission reductions; the total release averted by these interventions was between 0.6 and 1.8 billion metric tonnes. The study focused on interventions with at least 4.5% emission reductions, but the researchers noted that meeting the reductions required by the Paris Agreement would require 23 billion metric tonnes per year. Generally, carbon pricing was found to be most effective in developed countries, while regulation was most effective in the developing countries. Complementary policy mixes benefited from synergies, and were mostly found to be more effective interventions than the implementation of isolated policies.
The OECD recognise 48 distinct climate mitigation policies suitable for implementation at national level. Broadly, these can be categorised into three types: market based instruments, non market based instruments and other policies.
Other policies include the Establishing an Independent climate advisory body.
Non market based policies include the Implementing or tighening of Regulatory standards. These set technology or performance standards. They can be effective in addressing the market failure of informational barriers.
Among market based policies, the carbon price has been found to be the most effective (at least for developed economies), and has its own section below. Additional market based policy instruments for climate change mitigation include:
Emissions taxes These often require domestic emitters to pay a fixed fee or tax for every tonne of CO2 emissions they release into the atmosphere. Methane emissions from fossil fuel extraction are also occasionally taxed. But methane and nitrous oxide from agriculture are typically not subject to tax.
Removing unhelpful subsidies: Many countries provide subsidies for activities that affect emissions. For example, significant fossil fuel subsidies are present in many countries. Phasing-out fossil fuel subsidies is crucial to address the climate crisis. It must however be done carefully to avoid protests and making poor people poorer.
Creating helpful subsidies: Creating subsidies and financial incentives. One example is energy subsidies to support clean generation which is not yet commercially viable such as tidal power.
Tradable permits: A permit system can limit emissions.
Carbon pricing
Imposing additional costs on greenhouse gas emissions can make fossil fuels less competitive and accelerate investments into low-carbon sources of energy. A growing number of countries raise a fixed carbon tax or participate in dynamic carbon emission trading (ETS) systems. In 2021, more than 21% of global greenhouse gas emissions were covered by a carbon price. This was a big increase from earlier due to the introduction of the Chinese national carbon trading scheme.
Trading schemes offer the possibility to limit emission allowances to certain reduction targets. However, an oversupply of allowances keeps most ETS at low price levels around $10 with a low impact. This includes the Chinese ETS which started with $7/t in 2021. One exception is the European Union Emission Trading Scheme where prices began to rise in 2018. They reached about €80/t in 2022. This results in additional costs of about €0.04/KWh for coal and €0.02/KWh for gas combustion for electricity, depending on the emission intensity. Industries which have high energy requirements and high emissions often pay only very low energy taxes, or even none at all.
While this is often part of national schemes, carbon offsets and credits can be part of a voluntary market as well such as on the international market. Notably, the company Blue Carbon of the UAE has bought ownership over an area equivalent to the United Kingdom to be preserved in return for carbon credits.
International agreements
Almost all countries are parties to the United Nations Framework Convention on Climate Change (UNFCCC). The ultimate objective of the UNFCCC is to stabilize atmospheric concentrations of greenhouse gases at a level that would prevent dangerous human interference with the climate system.
Although not designed for this purpose, the Montreal Protocol has benefited climate change mitigation efforts. The Montreal Protocol is an international treaty that has successfully reduced emissions of ozone-depleting substances such as CFCs. These are also greenhouse gases.
Paris Agreement
History
Historically efforts to deal with climate change have taken place at a multinational level. They involve attempts to reach a consensus decision at the United Nations, under the United Nations Framework Convention on Climate Change (UNFCCC). This is the dominant approach historically of engaging as many international governments as possible in taking action on a worldwide public issue. The Montreal Protocol in 1987 is a precedent that this approach can work. But some critics say the top-down framework of only utilizing the UNFCCC consensus approach is ineffective. They put forward counter-proposals of bottom-up governance. At this same time this would lessen the emphasis on the UNFCCC.
The Kyoto Protocol to the UNFCCC adopted in 1997 set out legally binding emission reduction commitments for the "Annex 1" countries. The Protocol defined three international policy instruments ("Flexibility Mechanisms") which could be used by the Annex 1 countries to meet their emission reduction commitments. According to Bashmakov, use of these instruments could significantly reduce the costs for Annex 1 countries in meeting their emission reduction commitments.
The Paris Agreement reached in 2015 succeeded the Kyoto Protocol which expired in 2020. Countries that ratified the Kyoto protocol committed to reduce their emissions of carbon dioxide and five other greenhouse gases, or engage in carbon emissions trading if they maintain or increase emissions of these gases.
In 2015, the UNFCCC's "structured expert dialogue" came to the conclusion that, "in some regions and vulnerable ecosystems, high risks are projected even for warming above 1.5 °C". Together with the strong diplomatic voice of the poorest countries and the island nations in the Pacific, this expert finding was the driving force leading to the decision of the 2015 Paris Climate Conference to lay down this 1.5 °C long-term target on top of the existing 2 °C goal.
Society and culture
Commitments to divest
More than 1000 organizations with investments worth US$8 trillion have made commitments to fossil fuel divestment. Socially responsible investing funds allow investors to invest in funds that meet high environmental, social and corporate governance (ESG) standards.
Barriers
There are individual, institutional and market barriers to achieving climate change mitigation. They differ for all the different mitigation options, regions and societies.
Difficulties with accounting for carbon dioxide removal can act as economic barriers. This would apply to BECCS (bioenergy with carbon capture and storage). The strategies that companies follow can act as a barrier. But they can also accelerate decarbonisation.
In order to decarbonise societies the state needs to play a predominant role. This is because it requires a massive coordination effort. This strong government role can only work well if there is social cohesion, political stability and trust.
For land-based mitigation options, finance is a major barrier. Other barriers are cultural values, governance, accountability and institutional capacity.
Developing countries face further barriers to mitigation.
The cost of capital increased in the early 2020s. A lack of available capital and finance is common in developing countries. Together with the absence of regulatory standards, this barrier supports the proliferation of inefficient equipment.
There are also financial and capacity barrier in many of these countries.
One study estimates that only 0.12% of all funding for climate-related research goes on the social science of climate change mitigation. Vastly more funding goes on natural science studies of climate change. Considerable sums also go on studies of the impact of climate change and adaptation to it.
Impacts of the COVID-19 pandemic
The COVID-19 pandemic led some governments to shift their focus away from climate action, at least temporarily. This obstacle to environmental policy efforts may have contributed to slowed investment in green energy technologies. The economic slowdown resulting from COVID-19 added to this effect.
In 2020, carbon dioxide emissions fell by 6.4% or 2.3 billion tonnes globally. Greenhouse gas emissions rebounded later in the pandemic as many countries began lifting restrictions. The direct impact of pandemic policies had a negligible long-term impact on climate change.
Examples by country
United States
China
China has committed to peak emissions by 2030 and reach net zero by 2060. Warming cannot be limited to 1.5 °C if any coal plants in China (without carbon capture) operate after 2045. The Chinese national carbon trading scheme started in 2021.
European Union
The European Commission estimates that an additional €477 million in annual investment is needed for the European Union to meet its Fit-for-55 decarbonization goals.
In the European Union, government-driven policies and the European Green Deal have helped position greentech (as an example) as a vital area for venture capital investment. By 2023, venture capital in the EU's greentech sector equaled that of the United States, reflecting a concerted effort to drive innovation and mitigate climate change through targeted financial support. The European Green Deal has fostered policies that contributed to a 30% rise in venture capital for greentech companies in the EU from 2021 to 2023, despite a downturn in other sectors during the same period.
While overall venture capital investment in the EU remains about six times lower than in the United States, the greentech sector has closed this gap significantly, attracting substantial funding. Key areas benefitting from increased investments are energy storage, circular economy initiatives, and agricultural technology. This is supported by the EU's ambitious goal to reduce greenhouse gas emissions by at least 55% by 2030.
Related approaches
Relationship with solar radiation modification (SRM)
While solar radiation modification (SRM) could reduce surface temperatures, it temporarily masks climate change rather than addressing the root cause, which is greenhouse gases. SRM would work by altering how much solar radiation the Earth absorbs. Examples include reducing the amount of sunlight reaching the surface, reducing the optical thickness and lifetime of clouds, and changing the ability of the surface to reflect radiation. The IPCC describes SRM as a climate risk reduction strategy or supplementary option rather than a climate mitigation option.
The terminology in this area is still evolving. Experts sometimes use the term geoengineering or climate engineering in the scientific literature for both CDR or SRM, if the techniques are used at a global scale. IPCC reports no longer use the terms geoengineering or climate engineering.
| Physical sciences | Climate change | Earth science |
2119219 | https://en.wikipedia.org/wiki/Differentiable%20manifold | Differentiable manifold | In mathematics, a differentiable manifold (also differential manifold) is a type of manifold that is locally similar enough to a vector space to allow one to apply calculus. Any manifold can be described by a collection of charts (atlas). One may then apply ideas from calculus while working within the individual charts, since each chart lies within a vector space to which the usual rules of calculus apply. If the charts are suitably compatible (namely, the transition from one chart to another is differentiable), then computations done in one chart are valid in any other differentiable chart.
In formal terms, a differentiable manifold is a topological manifold with a globally defined differential structure. Any topological manifold can be given a differential structure locally by using the homeomorphisms in its atlas and the standard differential structure on a vector space. To induce a global differential structure on the local coordinate systems induced by the homeomorphisms, their compositions on chart intersections in the atlas must be differentiable functions on the corresponding vector space. In other words, where the domains of charts overlap, the coordinates defined by each chart are required to be differentiable with respect to the coordinates defined by every chart in the atlas. The maps that relate the coordinates defined by the various charts to one another are called transition maps.
The ability to define such a local differential structure on an abstract space allows one to extend the definition of differentiability to spaces without global coordinate systems. A locally differential structure allows one to define the globally differentiable tangent space, differentiable functions, and differentiable tensor and vector fields.
Differentiable manifolds are very important in physics. Special kinds of differentiable manifolds form the basis for physical theories such as classical mechanics, general relativity, and Yang–Mills theory. It is possible to develop a calculus for differentiable manifolds. This leads to such mathematical machinery as the exterior calculus. The study of calculus on differentiable manifolds is known as differential geometry.
"Differentiability" of a manifold has been given several meanings, including: continuously differentiable, k-times differentiable, smooth (which itself has many meanings), and analytic.
History
The emergence of differential geometry as a distinct discipline is generally credited to Carl Friedrich Gauss and Bernhard Riemann. Riemann first described manifolds in his famous habilitation lecture before the faculty at Göttingen. He motivated the idea of a manifold by an intuitive process of varying a given object in a new direction, and presciently described the role of coordinate systems and charts in subsequent formal developments:
Having constructed the notion of a manifoldness of n dimensions, and found that its true character consists in the property that the determination of position in it may be reduced to n determinations of magnitude, ... – B. Riemann
The works of physicists such as James Clerk Maxwell, and mathematicians Gregorio Ricci-Curbastro and Tullio Levi-Civita led to the development of tensor analysis and the notion of covariance, which identifies an intrinsic geometric property as one that is invariant with respect to coordinate transformations. These ideas found a key application in Albert Einstein's theory of general relativity and its underlying equivalence principle. A modern definition of a 2-dimensional manifold was given by Hermann Weyl in his 1913 book on Riemann surfaces. The widely accepted general definition of a manifold in terms of an atlas is due to Hassler Whitney.
Definition
Atlases
Let be a topological space. A chart on consists of an open subset of , and a homeomorphism from to an open subset of some Euclidean space . Somewhat informally, one may refer to a chart , meaning that the image of is an open subset of , and that is a homeomorphism onto its image; in the usage of some authors, this may instead mean that is itself a homeomorphism.
The presence of a chart suggests the possibility of doing differential calculus on ; for instance, if given a function and a chart on , one could consider the composition , which is a real-valued function whose domain is an open subset of a Euclidean space; as such, if it happens to be differentiable, one could consider its partial derivatives.
This situation is not fully satisfactory for the following reason. Consider a second chart on , and suppose that and contain some points in common. The two corresponding functions and are linked in the sense that they can be reparametrized into one another:
the natural domain of the right-hand side being . Since and are homeomorphisms, it follows that is a homeomorphism from to . Consequently it's just a bicontinuous function, thus even if both functions and are differentiable, their differential properties will not necessarily be strongly linked to one another, as is not guaranteed to be sufficiently differentiable for being able to compute the partial derivatives of the LHS applying the chain rule to the RHS. The same problem is found if one considers instead functions ; one is led to the reparametrization formula
at which point one can make the same observation as before.
This is resolved by the introduction of a "differentiable atlas" of charts, which specifies a collection of charts on for which the transition maps are all differentiable. This makes the situation quite clean: if is differentiable, then due to the first reparametrization formula listed above, the map is also differentiable on the region , and vice versa. Moreover, the derivatives of these two maps are linked to one another by the chain rule. Relative to the given atlas, this facilitates a notion of differentiable mappings whose domain or range is , as well as a notion of the derivative of such maps.
Formally, the word "differentiable" is somewhat ambiguous, as it is taken to mean different things by different authors; sometimes it means the existence of first derivatives, sometimes the existence of continuous first derivatives, and sometimes the existence of infinitely many derivatives. The following gives a formal definition of various (nonambiguous) meanings of "differentiable atlas". Generally, "differentiable" will be used as a catch-all term including all of these possibilities, provided .
Since every real-analytic map is smooth, and every smooth map is for any , one can see that any analytic atlas can also be viewed as a smooth atlas, and every smooth atlas can be viewed as a atlas. This chain can be extended to include holomorphic atlases, with the understanding that any holomorphic map between open subsets of can be viewed as a real-analytic map between open subsets of .
Given a differentiable atlas on a topological space, one says that a chart is differentiably compatible with the atlas, or differentiable relative to the given atlas, if the inclusion of the chart into the collection of charts comprising the given differentiable atlas results in a differentiable atlas. A differentiable atlas determines a maximal differentiable atlas, consisting of all charts which are differentiably compatible with the given atlas. A maximal atlas is always very large. For instance, given any chart in a maximal atlas, its restriction to an arbitrary open subset of its domain will also be contained in the maximal atlas. A maximal smooth atlas is also known as a smooth structure; a maximal holomorphic atlas is also known as a complex structure.
An alternative but equivalent definition, avoiding the direct use of maximal atlases, is to consider equivalence classes of differentiable atlases, in which two differentiable atlases are considered equivalent if every chart of one atlas is differentiably compatible with the other atlas. Informally, what this means is that in dealing with a smooth manifold, one can work with a single differentiable atlas, consisting of only a few charts, with the implicit understanding that many other charts and differentiable atlases are equally legitimate.
According to the invariance of domain, each connected component of a topological space which has a differentiable atlas has a well-defined dimension . This causes a small ambiguity in the case of a holomorphic atlas, since the corresponding dimension will be one-half of the value of its dimension when considered as an analytic, smooth, or atlas. For this reason, one refers separately to the "real" and "complex" dimension of a topological space with a holomorphic atlas.
Manifolds
A differentiable manifold is a Hausdorff and second countable topological space , together with a maximal differentiable atlas on . Much of the basic theory can be developed without the need for the Hausdorff and second countability conditions, although they are vital for much of the advanced theory. They are essentially equivalent to the general existence of bump functions and partitions of unity, both of which are used ubiquitously.
The notion of a manifold is identical to that of a topological manifold. However, there is a notable distinction to be made. Given a topological space, it is meaningful to ask whether or not it is a topological manifold. By contrast, it is not meaningful to ask whether or not a given topological space is (for instance) a smooth manifold, since the notion of a smooth manifold requires the specification of a smooth atlas, which is an additional structure. It could, however, be meaningful to say that a certain topological space cannot be given the structure of a smooth manifold. It is possible to reformulate the definitions so that this sort of imbalance is not present; one can start with a set (rather than a topological space ), using the natural analogue of a smooth atlas in this setting to define the structure of a topological space on .
Patching together Euclidean pieces to form a manifold
One can reverse-engineer the above definitions to obtain one perspective on the construction of manifolds. The idea is to start with the images of the charts and the transition maps, and to construct the manifold purely from this data. As in the above discussion, we use the "smooth" context but everything works just as well in other settings.
Given an indexing set let be a collection of open subsets of and for each let be an open (possibly empty) subset of and let be a smooth map. Suppose that is the identity map, that is the identity map, and that is the identity map. Then define an equivalence relation on the disjoint union by declaring to be equivalent to With some technical work, one can show that the set of equivalence classes can naturally be given a topological structure, and that the charts used in doing so form a smooth atlas. For the patching together the analytic structures(subset), see analytic varieties.
Differentiable functions
A real valued function f on an n-dimensional differentiable manifold M is called differentiable at a point if it is differentiable in any coordinate chart defined around p. In more precise terms, if is a differentiable chart where is an open set in containing p and is the map defining the chart, then f is differentiable at p if and only if
is differentiable at , that is is a differentiable function from the open set , considered as a subset of , to . In general, there will be many available charts; however, the definition of differentiability does not depend on the choice of chart at p. It follows from the chain rule applied to the transition functions between one chart and another that if f is differentiable in any particular chart at p, then it is differentiable in all charts at p. Analogous considerations apply to defining Ck functions, smooth functions, and analytic functions.
Differentiation of functions
There are various ways to define the derivative of a function on a differentiable manifold, the most fundamental of which is the directional derivative. The definition of the directional derivative is complicated by the fact that a manifold will lack a suitable affine structure with which to define vectors. Therefore, the directional derivative looks at curves in the manifold instead of vectors.
Directional differentiation
Given a real valued function f on an n dimensional differentiable manifold M, the directional derivative of f at a point p in M is defined as follows. Suppose that γ(t) is a curve in M with , which is differentiable in the sense that its composition with any chart is a differentiable curve in Rn. Then the directional derivative of f at p along γ is
If γ1 and γ2 are two curves such that , and in any coordinate chart ,
then, by the chain rule, f has the same directional derivative at p along γ1 as along γ2. This means that the directional derivative depends only on the tangent vector of the curve at p. Thus, the more abstract definition of directional differentiation adapted to the case of differentiable manifolds ultimately captures the intuitive features of directional differentiation in an affine space.
Tangent vector and the differential
A tangent vector at is an equivalence class of differentiable curves γ with , modulo the equivalence relation of first-order contact between the curves. Therefore,
in every coordinate chart . Therefore, the equivalence classes are curves through p with a prescribed velocity vector at p. The collection of all tangent vectors at p forms a vector space: the tangent space to M at p, denoted TpM.
If X is a tangent vector at p and f a differentiable function defined near p, then differentiating f along any curve in the equivalence class defining X gives a well-defined directional derivative along X:
Once again, the chain rule establishes that this is independent of the freedom in selecting γ from the equivalence class, since any curve with the same first order contact will yield the same directional derivative.
If the function f is fixed, then the mapping
is a linear functional on the tangent space. This linear functional is often denoted by df(p) and is called the differential of f at p:
Definition of tangent space and differentiation in local coordinates
Let be a topological -manifold with a smooth atlas Given let denote A "tangent vector at " is a mapping here denoted such that
for all Let the collection of tangent vectors at be denoted by Given a smooth function , define by sending a tangent vector to the number given by
which due to the chain rule and the constraint in the definition of a tangent vector does not depend on the choice of
One can check that naturally has the structure of a -dimensional real vector space, and that with this structure, is a linear map. The key observation is that, due to the constraint appearing in the definition of a tangent vector, the value of for a single element of automatically determines for all
The above formal definitions correspond precisely to a more informal notation which appears often in textbooks, specifically
and
With the idea of the formal definitions understood, this shorthand notation is, for most purposes, much easier to work with.
Partitions of unity
One of the topological features of the sheaf of differentiable functions on a differentiable manifold is that it admits partitions of unity. This distinguishes the differential structure on a manifold from stronger structures (such as analytic and holomorphic structures) that in general fail to have partitions of unity.
Suppose that M is a manifold of class Ck, where . Let {Uα} be an open covering of M. Then a partition of unity subordinate to the cover {Uα} is a collection of real-valued Ck functions φi on M satisfying the following conditions:
The supports of the φi are compact and locally finite;
The support of φi is completely contained in Uα for some α;
The φi sum to one at each point of M:
(Note that this last condition is actually a finite sum at each point because of the local finiteness of the supports of the φi.)
Every open covering of a Ck manifold M has a Ck partition of unity. This allows for certain constructions from the topology of Ck functions on Rn to be carried over to the category of differentiable manifolds. In particular, it is possible to discuss integration by choosing a partition of unity subordinate to a particular coordinate atlas, and carrying out the integration in each chart of Rn. Partitions of unity therefore allow for certain other kinds of function spaces to be considered: for instance Lp spaces, Sobolev spaces, and other kinds of spaces that require integration.
Differentiability of mappings between manifolds
Suppose M and N are two differentiable manifolds with dimensions m and n, respectively, and f is a function from M to N. Since differentiable manifolds are topological spaces we know what it means for f to be continuous. But what does "f is " mean for ? We know what that means when f is a function between Euclidean spaces, so if we compose f with a chart of M and a chart of N such that we get a map that goes from Euclidean space to M to N to Euclidean space we know what it means for that map to be . We define "f is " to mean that all such compositions of f with charts are . Once again, the chain rule guarantees that the idea of differentiability does not depend on which charts of the atlases on M and N are selected. However, defining the derivative itself is more subtle. If M or N is itself already a Euclidean space, then we don't need a chart to map it to one.
Bundles
Tangent bundle
The tangent space of a point consists of the possible directional derivatives at that point, and has the same dimension n as does the manifold. For a set of (non-singular) coordinates xk local to the point, the coordinate derivatives define a holonomic basis of the tangent space. The collection of tangent spaces at all points can in turn be made into a manifold, the tangent bundle, whose dimension is 2n. The tangent bundle is where tangent vectors lie, and is itself a differentiable manifold. The Lagrangian is a function on the tangent bundle. One can also define the tangent bundle as the bundle of 1-jets from R (the real line) to M.
One may construct an atlas for the tangent bundle consisting of charts based on , where Uα denotes one of the charts in the atlas for M. Each of these new charts is the tangent bundle for the charts Uα. The transition maps on this atlas are defined from the transition maps on the original manifold, and retain the original differentiability class.
Cotangent bundle
The dual space of a vector space is the set of real valued linear functions on the vector space. The cotangent space at a point is the dual of the tangent space at that point and the elements are referred to as cotangent vectors; the cotangent bundle is the collection of all cotangent vectors, along with the natural differentiable manifold structure.
Like the tangent bundle, the cotangent bundle is again a differentiable manifold. The Hamiltonian is a scalar on the cotangent bundle. The total space of a cotangent bundle has the structure of a symplectic manifold. Cotangent vectors are sometimes called covectors. One can also define the cotangent bundle as the bundle of 1-jets of functions from M to R.
Elements of the cotangent space can be thought of as infinitesimal displacements: if f is a differentiable function we can define at each point p a cotangent vector dfp, which sends a tangent vector Xp to the derivative of f associated with Xp. However, not every covector field can be expressed this way. Those that can are referred to as exact differentials. For a given set of local coordinates xk, the differentials dx form a basis of the cotangent space at p.
Tensor bundle
The tensor bundle is the direct sum of all tensor products of the tangent bundle and the cotangent bundle. Each element of the bundle is a tensor field, which can act as a multilinear operator on vector fields, or on other tensor fields.
The tensor bundle is not a differentiable manifold in the traditional sense, since it is infinite dimensional. It is however an algebra over the ring of scalar functions. Each tensor is characterized by its ranks, which indicate how many tangent and cotangent factors it has. Sometimes these ranks are referred to as covariant and contravariant ranks, signifying tangent and cotangent ranks, respectively.
Frame bundle
A frame (or, in more precise terms, a tangent frame), is an ordered basis of particular tangent space. Likewise, a tangent frame is a linear isomorphism of Rn to this tangent space. A moving tangent frame is an ordered list of vector fields that give a basis at every point of their domain. One may also regard a moving frame as a section of the frame bundle F(M), a principal bundle made up of the set of all frames over M. The frame bundle is useful because tensor fields on M can be regarded as equivariant vector-valued functions on F(M).
Jet bundles
On a manifold that is sufficiently smooth, various kinds of jet bundles can also be considered. The (first-order) tangent bundle of a manifold is the collection of curves in the manifold modulo the equivalence relation of first-order contact. By analogy, the k-th order tangent bundle is the collection of curves modulo the relation of k-th order contact. Likewise, the cotangent bundle is the bundle of 1-jets of functions on the manifold: the k-jet bundle is the bundle of their k-jets. These and other examples of the general idea of jet bundles play a significant role in the study of differential operators on manifolds.
The notion of a frame also generalizes to the case of higher-order jets. Define a k-th order frame to be the k-jet of a diffeomorphism from Rn to M. The collection of all k-th order frames, Fk(M), is a principal Gk bundle over M, where Gk is the group of k-jets; i.e., the group made up of k-jets of diffeomorphisms of Rn that fix the origin. Note that is naturally isomorphic to G1, and a subgroup of every Gk, . In particular, a section of F2(M) gives the frame components of a connection on M. Thus, the quotient bundle is the bundle of symmetric linear connections over M.
Calculus on manifolds
Many of the techniques from multivariate calculus also apply, mutatis mutandis, to differentiable manifolds. One can define the directional derivative of a differentiable function along a tangent vector to the manifold, for instance, and this leads to a means of generalizing the total derivative of a function: the differential. From the perspective of calculus, the derivative of a function on a manifold behaves in much the same way as the ordinary derivative of a function defined on a Euclidean space, at least locally. For example, there are versions of the implicit and inverse function theorems for such functions.
There are, however, important differences in the calculus of vector fields (and tensor fields in general). In brief, the directional derivative of a vector field is not well-defined, or at least not defined in a straightforward manner. Several generalizations of the derivative of a vector field (or tensor field) do exist, and capture certain formal features of differentiation in Euclidean spaces. The chief among these are:
The Lie derivative, which is uniquely defined by the differential structure, but fails to satisfy some of the usual features of directional differentiation.
An affine connection, which is not uniquely defined, but generalizes in a more complete manner the features of ordinary directional differentiation. Because an affine connection is not unique, it is an additional piece of data that must be specified on the manifold.
Ideas from integral calculus also carry over to differential manifolds. These are naturally expressed in the language of exterior calculus and differential forms. The fundamental theorems of integral calculus in several variables—namely Green's theorem, the divergence theorem, and Stokes' theorem—generalize to a theorem (also called Stokes' theorem) relating the exterior derivative and integration over submanifolds.
Differential calculus of functions
Differentiable functions between two manifolds are needed in order to formulate suitable notions of submanifolds, and other related concepts. If is a differentiable function from a differentiable manifold M of dimension m to another differentiable manifold N of dimension n, then the differential of f is a mapping . It is also denoted by Tf and called the tangent map. At each point of M, this is a linear transformation from one tangent space to another:
The rank of f at p is the rank of this linear transformation.
Usually the rank of a function is a pointwise property. However, if the function has maximal rank, then the rank will remain constant in a neighborhood of a point. A differentiable function "usually" has maximal rank, in a precise sense given by Sard's theorem. Functions of maximal rank at a point are called immersions and submersions:
If , and has rank m at , then f is called an immersion at p. If f is an immersion at all points of M and is a homeomorphism onto its image, then f is an embedding. Embeddings formalize the notion of M being a submanifold of N. In general, an embedding is an immersion without self-intersections and other sorts of non-local topological irregularities.
If , and has rank n at , then f is called a submersion at p. The implicit function theorem states that if f is a submersion at p, then M is locally a product of N and Rm−n near p. In formal terms, there exist coordinates in a neighborhood of f(p) in N, and functions x1, ..., xm−n defined in a neighborhood of p in M such that is a system of local coordinates of M in a neighborhood of p. Submersions form the foundation of the theory of fibrations and fibre bundles.
Lie derivative
A Lie derivative, named after Sophus Lie, is a derivation on the algebra of tensor fields over a manifold M. The vector space of all Lie derivatives on M forms an infinite dimensional Lie algebra with respect to the Lie bracket defined by
The Lie derivatives are represented by vector fields, as infinitesimal generators of flows (active diffeomorphisms) on M. Looking at it the other way around, the group of diffeomorphisms of M has the associated Lie algebra structure, of Lie derivatives, in a way directly analogous to the Lie group theory.
Exterior calculus
The exterior calculus allows for a generalization of the gradient, divergence and curl operators.
The bundle of differential forms, at each point, consists of all totally antisymmetric multilinear maps on the tangent space at that point. It is naturally divided into n-forms for each n at most equal to the dimension of the manifold; an n-form is an n-variable form, also called a form of degree n. The 1-forms are the cotangent vectors, while the 0-forms are just scalar functions. In general, an n-form is a tensor with cotangent rank n and tangent rank 0. But not every such tensor is a form, as a form must be antisymmetric.
Exterior derivative
The exterior derivative is a linear operator on the graded vector space of all smooth differential forms on a smooth manifold . It is usually denoted by . More precisely, if , for the operator maps the space of -forms on into the space of -forms (if there are no non-zero -forms on so the map is identically zero on -forms).
For example, the exterior differential of a smooth function is given in local coordinates , with associated local co-frame by the formula :
The exterior differential satisfies the following identity, similar to a product rule with respect to the wedge product of forms:
The exterior derivative also satisfies the identity . That is, if is a -form then the -form is identically vanishing. A form such that is called closed, while a form such that for some other form is called exact. Another formulation of the identity is that an exact form is closed. This allows one to define de Rham cohomology of the manifold , where the th cohomology group is the quotient group of the closed forms on by the exact forms on .
Topology of differentiable manifolds
Relationship with topological manifolds
Suppose that is a topological -manifold.
If given any smooth atlas , it is easy to find a smooth atlas which defines a different smooth manifold structure on consider a homeomorphism which is not smooth relative to the given atlas; for instance, one can modify the identity map using a localized non-smooth bump. Then consider the new atlas which is easily verified as a smooth atlas. However, the charts in the new atlas are not smoothly compatible with the charts in the old atlas, since this would require that and are smooth for any and with these conditions being exactly the definition that both and are smooth, in contradiction to how was selected.
With this observation as motivation, one can define an equivalence relation on the space of smooth atlases on by declaring that smooth atlases and are equivalent if there is a homeomorphism such that is smoothly compatible with and such that is smoothly compatible with
More briefly, one could say that two smooth atlases are equivalent if there exists a diffeomorphism in which one smooth atlas is taken for the domain and the other smooth atlas is taken for the range.
Note that this equivalence relation is a refinement of the equivalence relation which defines a smooth manifold structure, as any two smoothly compatible atlases are also compatible in the present sense; one can take to be the identity map.
If the dimension of is 1, 2, or 3, then there exists a smooth structure on , and all distinct smooth structures are equivalent in the above sense. The situation is more complicated in higher dimensions, although it isn't fully understood.
Some topological manifolds admit no smooth structures, as was originally shown with a ten-dimensional example by . A major application of partial differential equations in differential geometry due to Simon Donaldson, in combination with results of Michael Freedman, shows that many simply-connected compact topological 4-manifolds do not admit smooth structures. A well-known particular example is the E8 manifold.
Some topological manifolds admit many smooth structures which are not equivalent in the sense given above. This was originally discovered by John Milnor in the form of the exotic 7-spheres.
Classification
Every one-dimensional connected smooth manifold is diffeomorphic to either or each with their standard smooth structures.
For a classification of smooth 2-manifolds, see surface. A particular result is that every two-dimensional connected compact smooth manifold is diffeomorphic to one of the following: or or The situation is more nontrivial if one considers complex-differentiable structure instead of smooth structure.
The situation in three dimensions is quite a bit more complicated, and known results are more indirect. A remarkable result, proved in 2002 by methods of partial differential equations, is the geometrization conjecture, stating loosely that any compact smooth 3-manifold can be split up into different parts, each of which admits Riemannian metrics which possess many symmetries. There are also various "recognition results" for geometrizable 3-manifolds, such as Mostow rigidity and Sela's algorithm for the isomorphism problem for hyperbolic groups.
The classification of n-manifolds for n greater than three is known to be impossible, even up to homotopy equivalence. Given any finitely presented group, one can construct a closed 4-manifold having that group as fundamental group. Since there is no algorithm to decide the isomorphism problem for finitely presented groups, there is no algorithm to decide whether two 4-manifolds have the same fundamental group. Since the previously described construction results in a class of 4-manifolds that are homeomorphic if and only if their groups are isomorphic, the homeomorphism problem for 4-manifolds is undecidable. In addition, since even recognizing the trivial group is undecidable, it is not even possible in general to decide whether a manifold has trivial fundamental group, i.e. is simply connected.
Simply connected 4-manifolds have been classified up to homeomorphism by Freedman using the intersection form and Kirby–Siebenmann invariant. Smooth 4-manifold theory is known to be much more complicated, as the exotic smooth structures on R4 demonstrate.
However, the situation becomes more tractable for simply connected smooth manifolds of dimension ≥ 5, where the h-cobordism theorem can be used to reduce the classification to a classification up to homotopy equivalence, and surgery theory can be applied. This has been carried out to provide an explicit classification of simply connected 5-manifolds by Dennis Barden.
Structures on smooth manifolds
(Pseudo-)Riemannian manifolds
A Riemannian manifold consists of a smooth manifold together with a positive-definite inner product on each of the individual tangent spaces. This collection of inner products is called the Riemannian metric, and is naturally a symmetric 2-tensor field. This "metric" identifies a natural vector space isomorphism for each On a Riemannian manifold one can define notions of length, volume, and angle. Any smooth manifold can be given many different Riemannian metrics.
A pseudo-Riemannian manifold (also called a semi-Riemannian manifold) is a generalization of the notion of Riemannian manifold where the inner products are allowed to have an indefinite signature, as opposed to being positive-definite; they are still required to be non-degenerate. Every smooth pseudo-Riemannian and Riemmannian manifold defines a number of associated tensor fields, such as the Riemann curvature tensor. Lorentzian manifolds are pseudo-Riemannian manifolds of signature ; the case is fundamental in general relativity. Not every smooth manifold can be given a non-Riemannian pseudo-Riemannian structure; there are topological restrictions on doing so.
A Finsler manifold is a different generalization of a Riemannian manifold, in which the inner product is replaced with a vector norm; as such, this allows the definition of length, but not angle.
Symplectic manifolds
A symplectic manifold is a manifold equipped with a closed, nondegenerate 2-form. This condition forces symplectic manifolds to be even-dimensional, due to the fact that skew-symmetric matrices all have zero determinant. There are two basic examples:
Cotangent bundles, which arise as phase spaces in Hamiltonian mechanics, are a motivating example, since they admit a natural symplectic form.
All oriented two-dimensional Riemannian manifolds are, in a natural way, symplectic, by defining the form where, for any denotes the vector such that is an oriented -orthonormal basis of
Lie groups
A Lie group consists of a C∞ manifold together with a group structure on such that the product and inversion maps and are smooth as maps of manifolds. These objects often arise naturally in describing (continuous) symmetries, and they form an important source of examples of smooth manifolds.
Many otherwise familiar examples of smooth manifolds, however, cannot be given a Lie group structure, since given a Lie group and any , one could consider the map which sends the identity element to and hence, by considering the differential gives a natural identification between any two tangent spaces of a Lie group. In particular, by considering an arbitrary nonzero vector in one can use these identifications to give a smooth non-vanishing vector field on This shows, for instance, that no even-dimensional sphere can support a Lie group structure. The same argument shows, more generally, that every Lie group must be parallelizable.
Alternative definitions
Pseudogroups
The notion of a pseudogroup provides a flexible generalization of atlases in order to allow a variety of different structures to be defined on manifolds in a uniform way. A pseudogroup consists of a topological space S and a collection Γ consisting of homeomorphisms from open subsets of S to other open subsets of S such that
If , and U is an open subset of the domain of f, then the restriction f|U is also in Γ.
If f is a homeomorphism from a union of open subsets of S, , to an open subset of S, then provided for every i.
For every open , the identity transformation of U is in Γ.
If , then .
The composition of two elements of Γ is in Γ.
These last three conditions are analogous to the definition of a group. Note that Γ need not be a group, however, since the functions are not globally defined on S. For example, the collection of all local Ck diffeomorphisms on Rn form a pseudogroup. All biholomorphisms between open sets in Cn form a pseudogroup. More examples include: orientation preserving maps of Rn, symplectomorphisms, Möbius transformations, affine transformations, and so on. Thus, a wide variety of function classes determine pseudogroups.
An atlas (Ui, φi) of homeomorphisms φi from to open subsets of a topological space S is said to be compatible with a pseudogroup Γ provided that the transition functions are all in Γ.
A differentiable manifold is then an atlas compatible with the pseudogroup of Ck functions on Rn. A complex manifold is an atlas compatible with the biholomorphic functions on open sets in Cn. And so forth. Thus, pseudogroups provide a single framework in which to describe many structures on manifolds of importance to differential geometry and topology.
Structure sheaf
Sometimes, it can be useful to use an alternative approach to endow a manifold with a Ck-structure. Here k = 1, 2, ..., ∞, or ω for real analytic manifolds. Instead of considering coordinate charts, it is possible to start with functions defined on the manifold itself. The structure sheaf of M, denoted Ck, is a sort of functor that defines, for each open set , an algebra Ck(U) of continuous functions . A structure sheaf Ck is said to give M the structure of a Ck manifold of dimension n provided that, for any , there exists a neighborhood U of p and n functions such that the map is a homeomorphism onto an open set in Rn, and such that Ck|U is the pullback of the sheaf of k-times continuously differentiable functions on Rn.
In particular, this latter condition means that any function h in Ck(V), for V, can be written uniquely as , where H is a k-times differentiable function on f(V) (an open set in Rn). Thus, the sheaf-theoretic viewpoint is that the functions on a differentiable manifold can be expressed in local coordinates as differentiable functions on Rn, and a fortiori this is sufficient to characterize the differential structure on the manifold.
Sheaves of local rings
A similar, but more technical, approach to defining differentiable manifolds can be formulated using the notion of a ringed space. This approach is strongly influenced by the theory of schemes in algebraic geometry, but uses local rings of the germs of differentiable functions. It is especially popular in the context of complex manifolds.
We begin by describing the basic structure sheaf on Rn. If U is an open set in Rn, let
O(U) = Ck(U, R)
consist of all real-valued k-times continuously differentiable functions on U. As U varies, this determines a sheaf of rings on Rn. The stalk Op for consists of germs of functions near p, and is an algebra over R. In particular, this is a local ring whose unique maximal ideal consists of those functions that vanish at p. The pair is an example of a locally ringed space: it is a topological space equipped with a sheaf whose stalks are each local rings.
A differentiable manifold (of class Ck) consists of a pair where M is a second countable Hausdorff space, and OM is a sheaf of local R-algebras defined on M, such that the locally ringed space is locally isomorphic to . In this way, differentiable manifolds can be thought of as schemes modeled on Rn. This means that for each point , there is a neighborhood U of p, and a pair of functions , where
f : U → f(U) ⊂ Rn is a homeomorphism onto an open set in Rn.
f#: O|f(U) → f∗ (OM|U) is an isomorphism of sheaves.
The localization of f# is an isomorphism of local rings
f#f(p) : Of(p) → OM,p.
There are a number of important motivations for studying differentiable manifolds within this abstract framework. First, there is no a priori reason that the model space needs to be Rn. For example, (in particular in algebraic geometry), one could take this to be the space of complex numbers Cn equipped with the sheaf of holomorphic functions (thus arriving at the spaces of complex analytic geometry), or the sheaf of polynomials (thus arriving at the spaces of interest in complex algebraic geometry). In broader terms, this concept can be adapted for any suitable notion of a scheme (see topos theory). Second, coordinates are no longer explicitly necessary to the construction. The analog of a coordinate system is the pair , but these merely quantify the idea of local isomorphism rather than being central to the discussion (as in the case of charts and atlases). Third, the sheaf OM is not manifestly a sheaf of functions at all. Rather, it emerges as a sheaf of functions as a consequence of the construction (via the quotients of local rings by their maximal ideals). Hence, it is a more primitive definition of the structure (see synthetic differential geometry).
A final advantage of this approach is that it allows for natural direct descriptions of many of the fundamental objects of study to differential geometry and topology.
The cotangent space at a point is Ip/Ip2, where Ip is the maximal ideal of the stalk OM,p.
In general, the entire cotangent bundle can be obtained by a related technique (see cotangent bundle for details).
Taylor series (and jets) can be approached in a coordinate-independent manner using the Ip-adic filtration on OM,p.
The tangent bundle (or more precisely its sheaf of sections) can be identified with the sheaf of morphisms of OM into the ring of dual numbers.
Generalizations
The category of smooth manifolds with smooth maps lacks certain desirable properties, and people have tried to generalize smooth manifolds in order to rectify this. Diffeological spaces use a different notion of chart known as a "plot". Frölicher spaces and orbifolds are other attempts.
A rectifiable set generalizes the idea of a piece-wise smooth or rectifiable curve to higher dimensions; however, rectifiable sets are not in general manifolds.
Banach manifolds and Fréchet manifolds, in particular
manifolds of mappings
are infinite dimensional differentiable manifolds.
Non-commutative geometry
For a Ck manifold M, the set of real-valued Ck functions on the manifold forms an algebra under pointwise addition and multiplication, called the algebra of scalar fields or simply the algebra of scalars. This algebra has the constant function 1 as the multiplicative identity, and is a differentiable analog of the ring of regular functions in algebraic geometry.
It is possible to reconstruct a manifold from its algebra of scalars, first as a set, but also as a topological space – this is an application of the Banach–Stone theorem, and is more formally known as the spectrum of a C*-algebra. First, there is a one-to-one correspondence between the points of M and the algebra homomorphisms , as such a homomorphism φ corresponds to a codimension one ideal in Ck(M) (namely the kernel of φ), which is necessarily a maximal ideal. On the converse, every maximal ideal in this algebra is an ideal of functions vanishing at a single point, which demonstrates that MSpec (the Max Spec) of Ck(M) recovers M as a point set, though in fact it recovers M as a topological space.
One can define various geometric structures algebraically in terms of the algebra of scalars, and these definitions often generalize to algebraic geometry (interpreting rings geometrically) and operator theory (interpreting Banach spaces geometrically). For example, the tangent bundle to M can be defined as the derivations of the algebra of smooth functions on M.
This "algebraization" of a manifold (replacing a geometric object with an algebra) leads to the notion of a C*-algebra – a commutative C*-algebra being precisely the ring of scalars of a manifold, by Banach–Stone, and allows one to consider noncommutative C*-algebras as non-commutative generalizations of manifolds. This is the basis of the field of noncommutative geometry.
| Mathematics | Topology | null |
2119766 | https://en.wikipedia.org/wiki/Native%20copper | Native copper | Native copper is an uncombined form of copper that occurs as a natural mineral. Copper is one of the few metallic elements to occur in native form, although it most commonly occurs in oxidized states and mixed with other elements. Native copper was an important ore of copper in historic times and was used by pre-historic peoples.
Native copper occurs rarely as isometric cubic and octahedral crystals, but more typically as irregular masses and fracture fillings. It has a reddish, orangish, and/or brownish color on fresh surfaces, but typically is weathered and coated with a green tarnish of copper(II) carbonate (also known as patina or verdigris). Its specific gravity is 8.9 and its hardness is 2.5–3.
The mines of the Keweenaw native copper deposits of Upper Michigan were major copper producers in the 19th and early 20th centuries, and are the largest deposits of native copper in the world. Native Americans mined copper on a small scale at this and many other locations, and evidence exists of copper trading routes throughout North America among native peoples, proven by isotopic analysis. The first commercial mines in the Keweenaw Peninsula (which is nicknamed the "Copper Country" and "Copper Island") opened in the 1840s. Isle Royale in western Lake Superior was also a site of many tons of native copper. Some of it was extracted by native peoples, but only one of several commercial attempts at mining turned a profit there. An archived record of native copper originally found up river from Lake Superior, on the west branch of the Ontonagon River, via being dragged by a glacier is seen in the Ontonagon Boulder now in the possession of the Department of Mineral Sciences, National Museum of Natural History, Smithsonian Institution.
Another major native copper deposit is in Coro Coro, Bolivia.
The name copper comes from the Greek kyprios, "of Cyprus", the location of copper mines since pre-historic times.
Gallery
| Physical sciences | Minerals | Earth science |
2119871 | https://en.wikipedia.org/wiki/Emperor%20dragonfly | Emperor dragonfly | The emperor dragonfly or blue emperor (Anax imperator) is a large species of hawker dragonfly of the family Aeshnidae. It is the largest dragonfly in most of Europe, including the United Kingdom, although exceeded in some areas by other species.
Nomenclature
The generic name Anax is from the ancient Greek , "lord"; the specific epithet imperator is the Latin for "emperor", from imperare, to command.
Distribution
This dragonfly has a wide distribution through Afroeurasia; it is found throughout Africa and through most of Europe, the Arabian Peninsula, and south-western and central Asia. Since the 1990s, its range has expanded in Europe, both northwards and to higher altitudes. For example, the first Scandinavian record was in 1994 in Denmark; in 2002 it was first recorded in Sweden and in 2004 first in Scotland; today it is regular in all three countries. The species' northward expansion has been tied to global warming, and it is among the first odonata to do so.
Identification
The emperor dragonfly is a large and bulky species. It is long, with average being and males growing larger than females. The average wingspan is . When they first emerge, both sexes appear pale green with brown markings. The legs are brown with a yellow like base. Wings are born black but grow yellow-brown when they grow. Males have a bright sky blue or turquoise abdomen marked with a diagnostic black dorsal stripe. However, their blue colour may be faded during cold weather spells. The thorax and head of a male is apple green and their prominent eyes are blue. Females have similar markings but they are mainly a duller green. As the females age, their wings become browner. Less immediately visible features for both sexes are the yellow costa and brown spots on the wings. Emperor dragonflies can also be recognised by their flight patterns: they often fly with their abdomen hanging slightly downwards.
One of the largest species in Europe, the emperor dragonfly is exceeded by magnificent emperor, which occurs only marginally in the east Mediterranean and in length by females of the golden-ringed dragonfly, a species with an unusually long ovipositor. Thus, in most of Europe the emperor is the largest dragonfly species present.
Behaviour
They frequently fly high up into the sky in search of prey, which includes butterflies, other odonata and tadpoles. If their hunt is successful, they eat their smaller prey while flying. The dragonflies breed in a variety of aquatic habitats from large ponds to dikes and slow-moving rivers, but require a plentiful supply of vegetation in the water. They do sometimes breed in brackish water. The females lay the eggs into plants such as pondweed, and always lay alone. The aquatic larvae are very aggressive and are likely to influence the native species composition of freshwater ecosystems they arrive in. The larvae are also very large–around . The adult male is highly territorial, and difficult to approach.
Conservation
Emperor dragonflies are assessed as a least-concern species by the IUCN. The species has widespread and has a stable population.
Mitochrondrial genome
The mitogenome of the emperor dragonfly is the longest of all known dragonfly sequences. It has 16,087 base pairs. For comparison, the human mitogenome has 16,569 and the closely related dragonfly Anax parthenope has 15,366.
| Biology and health sciences | Odonata | Animals |
2121307 | https://en.wikipedia.org/wiki/Dirt%20road | Dirt road | A dirt road or track is a type of unpaved road not paved with asphalt, concrete, brick, or stone; made from the native material of the land surface through which it passes, known to highway engineers as subgrade material.
Terminology
Similar terms
Terms similar to dirt road are dry-weather road, earth road, or the "Class Four Highway" designation used in China. A track, dirt track, or earth track would normally be similar but less suitable for larger vehicles—the distinction is not well-defined. Laterite and murram roads, depending on material used, may be dirt roads or improved roads.
Improved road
Unpaved roads with a harder surface made by the addition of material such as gravel and aggregate (stones), might be referred to as dirt roads in common usage but are distinguished as improved roads by highway engineers. Improved unpaved roads include gravel roads and macadamized roads.
Characteristics
Compared to a gravel road, a dirt road is not usually graded regularly to produce an enhanced camber to encourage rainwater to drain off the road, and drainage ditches at the sides may be absent. They are unlikely to have embankments through low-lying areas. This leads to greater waterlogging and erosion, and after heavy rain the road may be impassable even to off-road vehicles. For this reason, in some countries, such as Australia and New Zealand and Finland, they are known as dry-weather roads.
Dirt roads take on different characteristics according to the soils and geology where they pass, and may be sandy, stony, rocky or have a bare earth surface, which could be extremely muddy and slippery when wet, and baked hard when dry. They are likely to become impassable after rain. They are common in rural areas of many countries, often very narrow and infrequently used, and are also found in metropolitan areas of many developing countries, where they may also be used as major highways and have considerable width.
Dirt roads almost always form a washboard-like surface with ridges. The reason for this is that dirt roads have tiny irregularities; a wheel hitting a bump pushes it forward, making it bigger, while a wheel pushing over a bump pushes dirt into the next bump. However, the surface can remain flat for velocities less than .
Driving on dirt roads
While most gravel roads are all-weather roads and can be used by ordinary cars, dirt roads may only be passable by trucks or four-wheel drive vehicles, especially in wet weather, or on rocky or very sandy sections. It is as easy to become bogged in sand as it is in mud; a high clearance under the vehicle may be required for rocky sections.
Driving on dirt roads requires great attention to variations in the surface and it is easier to lose control than on a gravel road.
Image gallery
| Technology | Road infrastructure | null |
2122837 | https://en.wikipedia.org/wiki/Ocimum%20tenuiflorum | Ocimum tenuiflorum | Ocimum tenuiflorum, commonly known as holy basil, tulsi or tulasi (), is an aromatic perennial plant in the family Lamiaceae. It is widely cultivated throughout the Southeast Asian tropics. It is native to tropical and subtropical regions of Asia, Australia and the western Pacific. This plant has escaped from cultivation and has naturalized in many tropical regions of the Americas. It is an agricultural and environmental weed.
Tulasi is cultivated for religious and traditional medicine purposes, and also for its essential oil. It is widely used as an herbal tea, commonly used in Ayurveda, and has a place within the Vaishnava tradition of Hinduism, in which devotees perform worship involving holy basil plants or leaves.
Morphology
Holy basil is an erect, many-branched subshrub, tall with hairy stems. Leaves are green or purple; they are simple, petioled, with an ovate blade up to long, which usually has a slightly toothed margin; they are strongly scented and have a decussate phyllotaxy. The purplish flowers are placed in close whorls on elongated racemes.
The three main morphotypes cultivated in India and Nepal are Ram tulsi (the most common type, with broad bright green leaves that are slightly sweet), the less common purplish green-leaved (Krishna or Shyam tulsi) and the common wild vana tulsi (e.g., Ocimum gratissimum).
Phytochemicals
The plant and its oil contain diverse phytochemicals, including tannins, flavonoids, eugenol, caryophyllenes, carvacrol, linalool, camphor, and cinnamyl acetate, among others. One study reported that the plant contains an eponymous family of 10 neolignan compounds called tulsinol A-J.
Specific aroma compounds in the essential oil are camphor (32%), eucalyptol (19%), ⍺-bisabolene (17%), eugenol (14%), germacrene (11%) and β-bisabolene (11%). In addition, more than 60 different aroma compounds were found through gas chromatography–mass spectrometry analysis of holy basil. However, other studies have stated tulsi essential oil consists mostly of eugenol (70%) β-elemene (11%), β-caryophyllene (8%), and germacrene (2%), with the balance being made up of various trace compounds, mostly terpenes.
Uses
Culinary
Tulasi has been used in Ayurvedic and Siddha practices for its supposed medicinal properties.
Thai cuisine
The leaves of holy basil, known as kaphrao in the Thai language (), are commonly used in Thai cuisine for certain stir-fries and curries such as phat kaphrao () — a stir-fry of Thai holy basil with meats, seafood or, as in khao phat kraphao, with rice. Two different types of holy basil are used in Thailand, a "red" variant which tends to be more pungent, and a "white" version for seafood dishes. Kaphrao should not be confused with horapha (), which is normally known as Thai basil, or with Thai lemon basil (maenglak; ).
Insect repellent
For centuries, the dried leaves have been mixed with stored grains to repel insects.
Nematicidal
The essential oil may have nematicidal properties against Tylenchulus semipenetrans, Meloidogyne javanica, Anguina tritici, and Heterodera cajani.
Disinfection
Water disinfection using O. tenuiflorum extracts was tested by Bhattacharjee et al. 2013 and Sadul et al. 2009. Both found an alcoholic extract to be more effective than aqueous or leaf juice. Sundaramurthi et al 2012 finds the result to be safe to drink and antimicrobial. A constituent analysis by Sadul found alkaloids, steroids, and tannins in the aqueous, and alkaloids and steroids only in the alcoholic extract.
Significance in Hinduism
Tulasi is a sacred plant for Hindus, particularly the Vaishnavite sect. It is worshipped as the avatar of Lakshmi, and is often planted in courtyards of Hindu houses or temples to Hanuman. The ritual lighting of lamps each evening during Kartik includes the worship of the tulsi plant. Vaishnavites are also known as "those who bear the tulsi around the neck".
Tulasi Vivaha is a ceremonial festival performed between Prabodhini Ekadashi (the 11th or 12th lunar day of the bright fortnight of the Hindu month of Kartika) and Kartik Purnima (the full moon of the month).
Every evening,Odia and Bengali Hindus place earthen lamps in front of tulsi plants. During the Kati Bihu festival celebrated in Assam, people light earthen lamps (diya) at the foot of the household tulsi plants and pray.
Gallery
| Biology and health sciences | Herbs and spices | Plants |
31716306 | https://en.wikipedia.org/wiki/Steam%20devil | Steam devil | A steam devil is a small, weak whirlwind over water (or sometimes wet land) that has drawn fog into the vortex, thus rendering it visible. They form over large lakes and oceans during cold air outbreaks while the water is still relatively warm, and can be an important mechanism in vertically transporting moisture. They are a component of sea smoke.
Smaller steam devils and steam whirls can form over geyser basins even in warm weather because of the very high water temperatures. Although observations of steam devils are generally quite rare, hot springs in Yellowstone Park produce them on a daily basis.
Steam devils have only been reported and studied since the 1970s. They are weaker than waterspouts and distinct from them. The latter are more akin to weak tornadoes over water.
Naming
Steam devils were first reported by Lyons and Pease in 1972 concerning their observations of Lake Michigan in January 1971. This month was a particularly cold one for Wisconsin (one of the coldest in the 20th century) which, combined with Lake Michigan staying mostly ice-free, produced good conditions for steam devil formation. Lyons and Pease named steam devils by comparison to the dust devils on land to which they have a comparable size and structure. They were also motivated by the need to distinguish steam devils from the much more powerful waterspout whose land equivalent is the tornado. Lyons and Pease wrote their article with the aim of persuading the National Oceanic and Atmospheric Administration to include steam devils in the International Field Year for the Great Lakes which was imminently to occur in 1972–3.
Appearance
Steam devils are vortices typically about 50 to 200 metres in diameter, essentially vertical, and up to 500 metres high. The general shape is like a small waterspout but they should not be considered related. Steam devils rotate with a cyclonic direction of motion, but not very fast or powerfully, usually just a few rotations per minute, and sometimes apparently not at all. There is usually a well-defined inner part of the rotating column of steam and a more ragged outer part from which clumps of steam often detach. Rather smaller steam devils can form over small lakes, especially the warm water in the hot springs of geyser basins. In these cases typical dimensions are a metre or so diameter, but can vary from less than 0.1 to 2 metres, and a height of 2 to 30 metres with a somewhat faster rotation of 60 rpm or so. The central core of the steam devil can be clear, in the same sense that the centre of a dust devil is clear of dust. The core is around 10% of the width of the rotating column. The sky above the steam devils may be clear, or there may be cumulus clouds present. In some cases the steam devils may rise directly into the cumulus, in these cases the cumulus may actually be caused by the steam devils - see below. Steam devils are a rare and short-lived phenomenon, typically surviving no more than three or four minutes, and the smaller ones over hot springs dissipating in a matter of seconds. Steam devils are sometimes confused with waterspouts as they can occur over the water.
Steam devils can become detached from their base and be blown downstream by the wind. On small bodies of water such as hot springs this can mean that the steam devil ends up over land away from the water altogether. Such steam devils continue to rotate even after they have become detached from the source of heat, but will soon dissipate.
Very small steam devils may have a poorly defined column and no identifiable clear inner core. Such vortices are more properly called steam whirls by analogy with the dust whirls of land.
Formation
A precondition for the formation of steam devils is the presence of a layer of moist air on the water with the misty air (called arctic steam fog) being drawn upwards into fog streamers (non-rotating columns of steam fog). For this to happen the body of water must be unfrozen, and thus relatively warm, and there must be some wind of cold, dry air to form the fog. The cold air is warmed by the water and is humidified by evaporation. The warmed air begins to rise, and as it does so is cooled adiabatically by the falling pressure causing the water vapour content to condense out into fog streamers.
For steam devils to form the air above the body of water must be very cold, and a fairly brisk (over 25 mph) wind of dry air needs to be blowing across the surface of the water. The temperature difference between the water and the air needs to be quite marked; the steam devils in figure 1 were forming with an air temperature of -21 °C (-6 °F) and a water temperature of 0.5 °C (33 °F) - a difference of 22 °C (39 °F). Under these conditions the air rises so energetically that the air flow becomes unstable and vortices start to form. Fog streamers drawn into the vortices render the vortices visible and they then become steam devils.
The steam fog tends to form irregular hexagonal cells in the horizontal plane which are elongated in the direction of the wind. In this honeycomb arrangement, three cells meet at a junction, and it is in these places that the steam devils form. This effect of vortex formation at the vertices of hexagonal cells is an example of vertex vortices.
The layer of cumulus seen above steam devils during cold air outbreaks on Lake Michigan and elsewhere may not be coincidental. Airborne radar studies during cold air outbreaks on the lake have shown that some steam devils penetrate through the thermal internal boundary layer (below which convective circulation takes place) and may be more significant for thermal mixing than normal convection, transporting moist air vertically above the convection boundary. The resulting large scale view is a layer of arctic steam fog close to the water surface, a layer of cumulus just above the convection boundary and a regular array of steam devils joining the two.
Occurrences
Steam devils are seen on the Great Lakes in early winter. They occur in the Atlantic off the coast of the Carolinas when cold air from the continent blows across the Gulf Stream. Steam devils can occur on small lakes and even over hot springs, but rather more rarely than on large bodies of water. It is also possible for steam devils to form over wet land if the air is cold and the sun is heating the ground.
Small steam devils occur at some of the larger hot springs in Yellowstone Park where a layer of steam fog hangs over the pools and wind can start to lift it up into fog streamers. One such example is the Grand Prismatic Spring in the Yellowstone Midway Geyser Basin. The air temperature can be high in terms of human comfort when the steam devils form. In 1982 a cluster of seventeen steam devils was observed when the air temperature was between 17 and 21 °C. Although this is much higher than, for instance, the temperature of the air over the Great Lakes, the water temperature is also proportionately higher, being very close to boiling, so the temperature difference is still 79 °C.
Another well known location in Yellowstone, the Old Faithful geyser, produces horizontal steam devils. In all, Yellowstone probably has the most frequent occurrences of accessible steam devils anywhere. Several steam devils are produced every hour at the most productive locations. Steam devils over geyser basins were first reported by Holle in 1977.
| Physical sciences | Storms | Earth science |
5318198 | https://en.wikipedia.org/wiki/Transmission%20coefficient | Transmission coefficient | The transmission coefficient is used in physics and electrical engineering when wave propagation in a medium containing discontinuities is considered. A transmission coefficient describes the amplitude, intensity, or total power of a transmitted wave relative to an incident wave.
Overview
Different fields of application have different definitions for the term. All the meanings are very similar in concept: In chemistry, the transmission coefficient refers to a chemical reaction overcoming a potential barrier; in optics and telecommunications it is the amplitude of a wave transmitted through a medium or conductor to that of the incident wave; in quantum mechanics it is used to describe the behavior of waves incident on a barrier, in a way similar to optics and telecommunications.
Although conceptually the same, the details in each field differ, and in some cases the terms are not an exact analogy.
Chemistry
In chemistry, in particular in transition state theory, there appears a certain "transmission coefficient" for overcoming a potential barrier. It is (often) taken to be unity for monomolecular reactions. It appears in the Eyring equation.
Optics
In optics, transmission is the property of a substance to permit the passage of light, with some or none of the incident light being absorbed in the process. If some light is absorbed by the substance, then the transmitted light will be a combination of the wavelengths of the light that was transmitted and not absorbed. For example, a blue light filter appears blue because it absorbs red and green wavelengths. If white light is shone through the filter, the light transmitted also appears blue because of the absorption of the red and green wavelengths.
The transmission coefficient is a measure of how much of an electromagnetic wave (light) passes through a surface or an optical element. Transmission coefficients can be calculated for either the amplitude or the intensity of the wave. Either is calculated by taking the ratio of the value after the surface or element to the value before. The transmission coefficient for total power is generally the same as the coefficient for intensity.
Telecommunications
In telecommunication, the transmission coefficient is the ratio of the amplitude of the complex transmitted wave to that of the incident wave at a discontinuity in the transmission line.
Consider a wave travelling through a transmission line with a step in impedance from to . When the wave transitions through the impedance step, a portion of the wave will be reflected back to the source. Because the voltage on a transmission line is always the sum of the forward and reflected waves at that point, if the incident wave amplitude is 1, and the reflected wave is , then the amplitude of the forward wave must be sum of the two waves or .
The value for is uniquely determined from first principles by noting that the incident power on the discontinuity must equal the sum of the power in the reflected and transmitted waves:
.
Solving the quadratic for leads both to the reflection coefficient:
,
and to the transmission coefficient:
.
The probability that a portion of a communications system, such as a line, circuit, channel or trunk, will meet specified performance criteria is also sometimes called the "transmission coefficient" of that portion of the system. The value of the transmission coefficient is inversely related to the quality of the line, circuit, channel or trunk.
Quantum mechanics
In non-relativistic quantum mechanics, the transmission coefficient and related reflection coefficient are used to describe the behavior of waves incident on a barrier. The transmission coefficient represents the probability flux of the transmitted wave relative to that of the incident wave. This coefficient is often used to describe the probability of a particle tunneling through a barrier.
The transmission coefficient is defined in terms of the incident and transmitted probability current density J according to:
where is the probability current in the wave incident upon the barrier with normal unit vector and is the probability current in the wave moving away from the barrier on the other side.
The reflection coefficient R is defined analogously:
Law of total probability requires that , which in one dimension reduces to the fact that the sum of the transmitted and reflected currents is equal in magnitude to the incident current.
For sample calculations, see rectangular potential barrier.
WKB approximation
Using the WKB approximation, one can obtain a tunnelling coefficient that looks like
where are the two classical turning points for the potential barrier. In the classical limit of all other physical parameters much larger than the reduced Planck constant, denoted , the transmission coefficient goes to zero. This classical limit would have failed in the situation of a square potential.
If the transmission coefficient is much less than 1, it can be approximated with the following formula:
where is the length of the barrier potential.
| Physical sciences | Optics | Physics |
5318699 | https://en.wikipedia.org/wiki/Fire%20lance | Fire lance | The fire lance () was a gunpowder weapon used by lighting it on fire, and is the ancestor of modern firearms. It first appeared in 10th–12th century China and was used to great effect during the Jin-Song Wars. It began as a small pyrotechnic device attached to a polearm weapon, used to gain a shock advantage at the start of a melee. As gunpowder improved, the explosive discharge was increased, and debris or pellets added, giving it some of the effects of a combination modern flamethrower and shotgun, but with a very short range (about ), and only one shot (although some were designed for two shots). By the late 13th century, fire lance barrels had transitioned to metal material to better withstand the explosive blast, and the lance-point was discarded in favor of relying solely on the gunpowder blast. These became the first hand cannons.
Design
The first fire lances consisted of a tube, usually bamboo, containing gunpowder and a slow match, strapped to a spear or other polearm weapon. Once ignited, the gunpowder tube would ideally eject a stream of flames in the direction of the spearhead. Projectiles such as iron pellets or pottery shards were later added to the gunpowder. Upon firing, the gunpowder charge ejected the projectiles along with the flame.
Metal fire lance barrels appeared around the mid-13th century and these began to be used independently of the lance itself. The independent metal barrel was known as an 'eruptor' and became the forerunner of the hand cannon.
In Europe, versions with wooden tubes were used.
History
China
The earliest evidence of fire lances appeared in China in the year 950. However usage of fire lances in warfare was not mentioned until 1132 when Song garrisons used them during the Siege of De'an, in modern-day Anlu, Hubei, in a sortie against the Jin dynasty (1115–1234).
In 1163, fire lances were attached to war carts known as "at-your-desire-carts" used to defend mobile firebomb trebuchets.
In the late 1100s, pieces of shrapnel such as porcelain shards and small iron pellets were added to the gunpowder tube. At some point fire lances discarded the spearhead altogether and relied solely on their firepower.
By 1232, the Jin were also using fire lances, but with improved reusable barrels consisting of durable paper material. According to the History of Jin, these fire lances had a range of roughly three meters:
In 1233, Jin soldiers used fire lances successfully against the Mongols. Pucha Guannu led 450 Jin fire lancers and routed an entire Mongol encampment. The Mongol soldiers were apparently disdainful of other Jin weapons, but greatly feared the fire lance.
In 1259, a pellet wad that occluded the barrel was recorded to have been used as a fire lance projectile, making it the first recorded bullet in history.
By 1276, fire lances had transitioned to metal barrels. Fire lances were also being used by cavalrymen at this point, as evidenced by the account of a Song-Yuan battle in which two fire lance armed Song cavalrymen rushed a Chinese officer of Bayan of the Baarin. The Huolongjing also mentions a gourd fire lance which was used by cavalrymen as well as foot soldiers.
The metal-barreled fire lance began to be used independently of the lance around the mid to late 13th century. These proto-cannons which fired co-viative projectiles, known as 'eruptors,' were the forerunners of the hand cannon.
Later history
By 1280, the Middle East had acquired fire lances.
In 1396, European knights took up fire lances as mounted weapons.
In 15th century Japanese samurai used fire lances.
The last recorded usage of fire lances in Europe occurred during the Storming of Bristol in 1643 although the Commonwealth of England was still issuing them to ships in 1660.
Troncks
Versions where the fireworks and shot were placed in a wooden tube at the end of a pole were known as Troncks, fire-trunks or bombas in Europe. The fireworks had alternating slow and fast burning sections.
They were frequently issued to warships and a surviving example was found in the wreck of the La Trinidad Valencera. Testing of an attempted reconstruction was carried out in 1988. During the test multiple sections of the Tronck ignited at once.
Gallery
| Technology | Firearms | null |
5321285 | https://en.wikipedia.org/wiki/Radium%20and%20radon%20in%20the%20environment | Radium and radon in the environment | Radium and radon are important contributors to environmental radioactivity. Radon occurs naturally as a result of decay of radioactive elements in soil and it can accumulate in houses built on areas where such decay occurs. Radon is a major cause of cancer; it is estimated to contribute to ~2% of all cancer related deaths in Europe.
Radium, like radon, is radioactive and is found in small quantities in nature and is hazardous to life if radiation exceeds 20-50 mSv/year. Radium is a decay product of uranium and thorium. Radium may also be released into the environment by human activity: for example, in improperly discarded products painted with radioluminescent paint.
Radium
In the oil and gas industries
Residues from the oil and gas industry often contain radium and its daughters. The sulfate scale from an oil well can be very radium rich. The water inside an oil field is often very rich in strontium, barium and radium, while seawater is very rich in sulfate: so if water from an oil well is discharged into the sea or mixed with seawater, the radium is likely to be brought out of solution by the barium/strontium sulfate which acts as a carrier precipitate.
Radioluminescent (glow in the dark) products
It is not unknown for local contamination to arise from improper disposal of radium-based radioluminescent paints.
In radioactive quackery
Eben Byers was a wealthy American socialite whose death in 1932 from using a radioactive quackery product called Radithor is a prominent example of a death caused by radium. Radithor contained ~1 μCi (40 kBq) of 226Ra and 1 μCi of 228Ra per bottle. Radithor was taken by mouth and radium, being a calcium mimic, has a very long biological halflife in bone.
Radon
Most of the dose is due to the decay of the polonium (218Po) and lead (214Pb) daughters of 222Rn. By controlling exposure to the daughters the radioactive dose to the skin and lungs can be reduced by at least 90%. This can be done by wearing a dust mask, and wearing a suit to cover the entire body. Note that exposure to smoke at the same time as radon and radon daughters will increase the harmful effect of the radon. In uranium miners radon has been found to be more carcinogenic in smokers than in non-smokers.
Occurrence
Radon concentration in open air varies between 1 and 100 Bq m−3. Radon can be found in some spring waters and hot springs. The towns of Misasa, Japan, and Bad Kreuznach, Germany boast radium-rich springs which emit radon, as does Radium Springs, New Mexico.
Radon exhausts naturally from the ground, particularly in certain regions, especially but not only regions with granitic soils. However, not all granitic regions are prone to high emissions of radon. For instance, while the rock which Aberdeen is on is very radium rich, the rock lacks the cracks required for the radon to migrate. In other nearby areas of Scotland (to the north of Aberdeen) and in Cornwall/Devon the radon is very much able to leave the rock.
Radon is a decay product of radium which in turn is a decay product of uranium. Maps of average radon levels in houses are available, to assist in planning mitigation measures.
While high uranium in the soil/rock under a house does not always lead to a high radon level in air, a positive correlation between the uranium content of the soil and the radon level in air can be seen.
In air
Radon harms indoor air quality in many homes. (See "In houses" below.)
Radon (222Rn) released into the air decays to 210Pb and other radioisotopes and the levels of 210Pb can be measured. It is important to note that the rate of deposition of this radioisotope is very dependent on the season. Here is a graph of the deposition rate observed in Japan.
In groundwater
Well water can be very rich in radon; the use of this water inside a house is another route allowing radon to enter the house. The radon can enter the air and then be a source of exposure to the humans, or the water can be consumed by humans which is a different exposure route.
Radon in rainwater
Rainwater can be highly radioactive due to high levels of radon and its decay progenies 214Bi and 214Pb; the concentrations of these radioisotopes can be high enough to seriously disrupt radiation monitoring at nuclear power plants. The highest levels of radon in rainwater occur during thunderstorms, and it is hypothesized that radon is concentrated in thunderstorms because it forms some positive ions during thunderstorms. Estimates of the age of raindrops have been obtained from measuring the isotopic abundance of radon's short-lived decay progeny in rainwater.
In the oil and gas industries
Water, oil and gas from a well often contain radon. The radon decays to form solid radioisotopes which form coatings on the inside of pipework. In an oil processing plant the area of the plant where propane is processed is often one of the more contaminated areas of the plant as radon has a similar boiling point to propane.
In mines
Because uranium minerals emit radon gas, and their harmful and highly radioactive decay products, uranium mining is considerably more dangerous than other (already dangerous) hard rock mining, requiring adequate ventilation systems if the mines are not open pit. In the 1950s, a significant number of American uranium miners were Navajo, as many uranium deposits were discovered on Navajo reservations. A statistically significant proportion of these miners later developed small-cell lung cancer, a type of cancer usually not associated with smoking, after exposure to uranium ore and radon-222, a natural decay product of uranium. The cancer causing agent has been shown to be the radon which is produced by the uranium, and not the uranium itself. Some survivors and their descendants received compensation under the Radiation Exposure Compensation Act in 1990.
The level of radon in the air of mines is now normally controlled by law. In a working mine, the radon level can be controlled by ventilation, sealing off old workings and controlling the water in the mine. The level in a mine can go up when a mine is abandoned; it can reach a level which can cause the skin to become red (a mild radiation burn). The radon levels in some of the mines can reach 400 to 700 kBq m−3.
A common unit of exposure of lung tissue to alpha emitters is the working level month (WLM), this is where the human lungs have been exposed for 170 hours (a typical month worth of work for a miner) to air which has 3.7 kBq of 222Rn (in equilibrium with its decay products). This is air which has the alpha dose rate of 1 working level (WL). It is estimated that the average person (general public) is subject to 0.2 WLM per year, which works out at about 15 to 20 WLM in a lifetime. According to the NRC, 1 WLM is a 5 to 10 mSv lung dose (0.5 to 1.0 rem), while the Organisation for Economic Co-operation and Development (OECD) consider that 1 WLM is equal to a lung dose of 5.5 mSv, and the International Commission on Radiological Protection (ICRP) consider 1 WLM to be a 5 mSv lung dose for professional workers (and a 4 mSv lung dose for the general public). Lastly the United Nations Scientific Committee on the Effects of Atomic Radiation (UNSCEAR) consider that the exposure of the lungs to 1 Bq of 222Rn (in equilibrium with its decay products) for one year will cause a dose of 61 μSv.
In houses
It has been known since at least the 1950s that radon is present in indoor air, and research into its effects on human health started in the early 1970s. The danger of radon exposure in dwellings received more widespread public awareness after 1984, as a result of the case of Stanley Watras, an employee at the Limerick nuclear power plant in Pennsylvania. Mr. Watras set off the radiation alarms (see Geiger counter) on his way into work for two weeks straight while authorities searched for the source of the contamination. They were shocked to find that the source was astonishingly high levels of radon in his basement and it was not related to the nuclear plant. The risks associated with living in his house were estimated to be equivalent to smoking 135 packs of cigarettes every day.
Depending how houses are built and ventilated, radon may accumulate in basements and dwellings. The European Union recommends that mitigation should be taken starting from concentrations of 400 Bq/m3 for old houses, and 200 Bq/m3 for new ones.
The National Council on Radiation Protection and Measurements (NCRP) recommends action for any house with a concentration higher than 8 pCi/L (300 Bq/m3).
The United States Environmental Protection Agency recommends action for any house with a concentration higher than 148 Bq/m3 (given as 4 pCi/L). Nearly one in 15 homes in the U.S. has a high level of indoor radon according to their statistics. The U.S. Surgeon General and EPA recommend all homes be tested for radon. Since 1985, millions of homes have been tested for radon in the U.S.
By adding a crawl space under the ground floor, which is subject to forced ventilation, the radon level in the house can be lowered.
| Physical sciences | Geochemistry | Earth science |
21647820 | https://en.wikipedia.org/wiki/Whole%20genome%20sequencing | Whole genome sequencing | Whole genome sequencing (WGS) is the process of determining the entirety, or nearly the entirety, of the DNA sequence of an organism's genome at a single time. This entails sequencing all of an organism's chromosomal DNA as well as DNA contained in the mitochondria and, for plants, in the chloroplast.
Whole genome sequencing has largely been used as a research tool, but was being introduced to clinics in 2014. In the future of personalized medicine, whole genome sequence data may be an important tool to guide therapeutic intervention. The tool of gene sequencing at SNP level is also used to pinpoint functional variants from association studies and improve the knowledge available to researchers interested in evolutionary biology, and hence may lay the foundation for predicting disease susceptibility and drug response.
Whole genome sequencing should not be confused with DNA profiling, which only determines the likelihood that genetic material came from a particular individual or group, and does not contain additional information on genetic relationships, origin or susceptibility to specific diseases. In addition, whole genome sequencing should not be confused with methods that sequence specific subsets of the genome – such methods include whole exome sequencing (1–2% of the genome) or SNP genotyping (< 0.1% of the genome).
History
The DNA sequencing methods used in the 1970s and 1980s were manual; for example, Maxam–Gilbert sequencing and Sanger sequencing. Several whole bacteriophage and animal viral genomes were sequenced by these techniques, but the shift to more rapid, automated sequencing methods in the 1990s facilitated the sequencing of the larger bacterial and eukaryotic genomes.
The first virus to have its complete genome sequenced was the Bacteriophage MS2 by 1976. In 1992, yeast chromosome III was the first chromosome of any organism to be fully sequenced. The first organism whose entire genome was fully sequenced was Haemophilus influenzae in 1995. After it, the genomes of other bacteria and some archaea were first sequenced, largely due to their small genome size. H. influenzae has a genome of 1,830,140 base pairs of DNA. In contrast, eukaryotes, both unicellular and multicellular such as Amoeba dubia and humans (Homo sapiens) respectively, have much larger genomes (see C-value paradox). Amoeba dubia has a genome of 700 billion nucleotide pairs spread across thousands of chromosomes. Humans contain fewer nucleotide pairs (about 3.2 billion in each germ cell – note the exact size of the human genome is still being revised) than A. dubia, however, their genome size far outweighs the genome size of individual bacteria.
The first bacterial and archaeal genomes, including that of H. influenzae, were sequenced by Shotgun sequencing. In 1996, the first eukaryotic genome (Saccharomyces cerevisiae) was sequenced. S. cerevisiae, a model organism in biology has a genome of only around 12 million nucleotide pairs, and was the first unicellular eukaryote to have its whole genome sequenced. The first multicellular eukaryote, and animal, to have its whole genome sequenced was the nematode worm: Caenorhabditis elegans in 1998. Eukaryotic genomes are sequenced by several methods including Shotgun sequencing of short DNA fragments and sequencing of larger DNA clones from DNA libraries such as bacterial artificial chromosomes (BACs) and yeast artificial chromosomes (YACs).
In 1999, the entire DNA sequence of human chromosome 22, the second shortest human autosome, was published. By the year 2000, the second animal and second invertebrate (yet first insect) genome was sequenced – that of the fruit fly Drosophila melanogaster – a popular choice of model organism in experimental research. The first plant genome – that of the model organism Arabidopsis thaliana – was also fully sequenced by 2000. By 2001, a draft of the entire human genome sequence was published. The genome of the laboratory mouse Mus musculus was completed in 2002.
In 2004, the Human Genome Project published an incomplete version of the human genome. In 2008, a group from Leiden, the Netherlands, reported the sequencing of the first female human genome (Marjolein Kriek).
Currently thousands of genomes have been wholly or partially sequenced.
Experimental details
Cells used for sequencing
Almost any biological sample containing a full copy of the DNA—even a very small amount of DNA or ancient DNA—can provide the genetic material necessary for full genome sequencing. Such samples may include saliva, epithelial cells, bone marrow, hair (as long as the hair contains a hair follicle), seeds, plant leaves, or anything else that has DNA-containing cells.
The genome sequence of a single cell selected from a mixed population of cells can be determined using techniques of single cell genome sequencing. This has important advantages in environmental microbiology in cases where a single cell of a particular microorganism species can be isolated from a mixed population by microscopy on the basis of its morphological or other distinguishing characteristics. In such cases the normally necessary steps of isolation and growth of the organism in culture may be omitted, thus allowing the sequencing of a much greater spectrum of organism genomes.
Single cell genome sequencing is being tested as a method of preimplantation genetic diagnosis, wherein a cell from the embryo created by in vitro fertilization is taken and analyzed before embryo transfer into the uterus. After implantation, cell-free fetal DNA can be taken by simple venipuncture from the mother and used for whole genome sequencing of the fetus.
Early techniques
Sequencing of nearly an entire human genome was first accomplished in 2000 partly through the use of shotgun sequencing technology. While full genome shotgun sequencing for small (4000–7000 base pair) genomes was already in use in 1979, broader application benefited from pairwise end sequencing, known colloquially as double-barrel shotgun sequencing. As sequencing projects began to take on longer and more complicated genomes, multiple groups began to realize that useful information could be obtained by sequencing both ends of a fragment of DNA. Although sequencing both ends of the same fragment and keeping track of the paired data was more cumbersome than sequencing a single end of two distinct fragments, the knowledge that the two sequences were oriented in opposite directions and were about the length of a fragment apart from each other was valuable in reconstructing the sequence of the original target fragment.
The first published description of the use of paired ends was in 1990 as part of the sequencing of the human HPRT locus, although the use of paired ends was limited to closing gaps after the application of a traditional shotgun sequencing approach. The first theoretical description of a pure pairwise end sequencing strategy, assuming fragments of constant length, was in 1991. In 1995, the innovation of using fragments of varying sizes was introduced, and demonstrated that a pure pairwise end-sequencing strategy would be possible on large targets. The strategy was subsequently adopted by The Institute for Genomic Research (TIGR) to sequence the entire genome of the bacterium Haemophilus influenzae in 1995, and then by Celera Genomics to sequence the entire fruit fly genome in 2000, and subsequently the entire human genome. Applied Biosystems, now called Life Technologies, manufactured the automated capillary sequencers utilized by both Celera Genomics and The Human Genome Project.
Current techniques
While capillary sequencing was the first approach to successfully sequence a nearly full human genome, it is still too expensive and takes too long for commercial purposes. Since 2005, capillary sequencing has been progressively displaced by high-throughput (formerly "next-generation") sequencing technologies such as Illumina dye sequencing, pyrosequencing, and SMRT sequencing. All of these technologies continue to employ the basic shotgun strategy, namely, parallelization and template generation via genome fragmentation.
Other technologies have emerged, including Nanopore technology. Though the sequencing accuracy of Nanopore technology is lower than those above, its read length is on average much longer. This generation of long reads is valuable especially in de novo whole-genome sequencing applications.
Analysis
In principle, full genome sequencing can provide the raw nucleotide sequence of an individual organism's DNA at a single point in time. However, further analysis must be performed to provide the biological or medical meaning of this sequence, such as how this knowledge can be used to help prevent disease. Methods for analyzing sequencing data are being developed and refined.
Because sequencing generates a lot of data (for example, there are approximately six billion base pairs in each human diploid genome), its output is stored electronically and requires a large amount of computing power and storage capacity.
While analysis of WGS data can be slow, it is possible to speed up this step by using dedicated hardware.
Commercialization
A number of public and private companies are competing to develop a full genome sequencing platform that is commercially robust for both research and clinical use, including Illumina, Knome, Sequenom,
454 Life Sciences, Pacific Biosciences, Complete Genomics,
Helicos Biosciences, GE Global Research (General Electric), Affymetrix, IBM, Intelligent Bio-Systems, Life Technologies, Oxford Nanopore Technologies, and the Beijing Genomics Institute. These companies are heavily financed and backed by venture capitalists, hedge funds, and investment banks.
A commonly-referenced commercial target for sequencing cost until the late 2010s was $1,000USD, however, the private companies are working to reach a new target of only $100.
Incentive
In October 2006, the X Prize Foundation, working in collaboration with the J. Craig Venter Science Foundation, established the Archon X Prize for Genomics, intending to award $10 million to "the first team that can build a device and use it to sequence 100 human genomes within 10 days or less, with an accuracy of no more than one error in every 1,000,000 bases sequenced, with sequences accurately covering at least 98% of the genome, and at a recurring cost of no more than $1,000 per genome". The Archon X Prize for Genomics was cancelled in 2013, before its official start date.
History
In 2007, Applied Biosystems started selling a new type of sequencer called SOLiD System. The technology allowed users to sequence 60 gigabases per run.
In June 2009, Illumina announced that they were launching their own Personal Full Genome Sequencing Service at a depth of 30× for $48,000 per genome. In August, the founder of Helicos Biosciences, Stephen Quake, stated that using the company's Single Molecule Sequencer he sequenced his own full genome for less than $50,000. In November, Complete Genomics published a peer-reviewed paper in Science demonstrating its ability to sequence a complete human genome for $1,700.
In May 2011, Illumina lowered its Full Genome Sequencing service to $5,000 per human genome, or $4,000 if ordering 50 or more.
Helicos Biosciences, Pacific Biosciences, Complete Genomics, Illumina, Sequenom, ION Torrent Systems, Halcyon Molecular, NABsys, IBM, and GE Global appear to all be going head to head in the race to commercialize full genome sequencing.
With sequencing costs declining, a number of companies began claiming that their equipment would soon achieve the $1,000 genome: these companies included Life Technologies in January 2012, Oxford Nanopore Technologies in February 2012, and Illumina in February 2014. In 2015, the NHGRI estimated the cost of obtaining a whole-genome sequence at around $1,500. In 2016, Veritas Genetics began selling whole genome sequencing, including a report as to some of the information in the sequencing for $999. In summer 2019, Veritas Genetics cut the cost for WGS to $599. In 2017, BGI began offering WGS for $600.
However, in 2015, some noted that effective use of whole gene sequencing can cost considerably more than $1000. Also, reportedly there remain parts of the human genome that have not been fully sequenced by 2017.
Comparison with other technologies
DNA microarrays
Full genome sequencing provides information on a genome that is orders of magnitude larger than by DNA arrays, the previous leader in genotyping technology.
For humans, DNA arrays currently provide genotypic information on up to one million genetic variants, while full genome sequencing will provide information on all six billion bases in the human genome, or 3,000 times more data. Because of this, full genome sequencing is considered a disruptive innovation to the DNA array markets as the accuracy of both range from 99.98% to 99.999% (in non-repetitive DNA regions) and their consumables cost of $5000 per 6 billion base pairs is competitive (for some applications) with DNA arrays ($500 per 1 million basepairs).
Applications
Mutation frequencies
Whole genome sequencing has established the mutation frequency for whole human genomes. The mutation frequency in the whole genome between generations for humans (parent to child) is about 70 new mutations per generation. An even lower level of variation was found comparing whole genome sequencing in blood cells for a pair of monozygotic (identical twins) 100-year-old centenarians. Only 8 somatic differences were found, though somatic variation occurring in less than 20% of blood cells would be undetected.
In the specifically protein coding regions of the human genome, it is estimated that there are about 0.35 mutations that would change the protein sequence between parent/child generations (less than one mutated protein per generation).
In cancer, mutation frequencies are much higher, due to genome instability. This frequency can further depend on patient age, exposure to DNA damaging agents (such as UV-irradiation or components of tobacco smoke) and the activity/inactivity of DNA repair mechanisms. Furthermore, mutation frequency can vary between cancer types: in germline cells, mutation rates occur at approximately 0.023 mutations per megabase, but this number is much higher in breast cancer (1.18-1.66 somatic mutations per Mb), in lung cancer (17.7) or in melanomas (≈33). Since the haploid human genome consists of approximately 3,200 megabases, this translates into about 74 mutations (mostly in noncoding regions) in germline DNA per generation, but 3,776-5,312 somatic mutations per haploid genome in breast cancer, 56,640 in lung cancer and 105,600 in melanomas.
The distribution of somatic mutations across the human genome is very uneven, such that the gene-rich, early-replicating regions receive fewer mutations than gene-poor, late-replicating heterochromatin, likely due to differential DNA repair activity. In particular, the histone modification H3K9me3 is associated with high, and H3K36me3 with low mutation frequencies.
Genome-wide association studies
In research, whole-genome sequencing can be used in a Genome-Wide Association Study (GWAS) – a project aiming to determine the genetic variant or variants associated with a disease or some other phenotype.
Diagnostic use
In 2009, Illumina released its first whole genome sequencers that were approved for clinical as opposed to research-only use and doctors at academic medical centers began quietly using them to try to diagnose what was wrong with people whom standard approaches had failed to help. In 2009, a team from Stanford led by Euan Ashley performed clinical interpretation of a full human genome, that of bioengineer Stephen Quake. In 2010, Ashley's team reported whole genome molecular autopsy and in 2011, extended the interpretation framework to a fully sequenced family, the West family, who were the first family to be sequenced on the Illumina platform. The price to sequence a genome at that time was $19,500USD, which was billed to the patient but usually paid for out of a research grant; one person at that time had applied for reimbursement from their insurance company. For example, one child had needed around 100 surgeries by the time he was three years old, and his doctor turned to whole genome sequencing to determine the problem; it took a team of around 30 people that included 12 bioinformatics experts, three sequencing technicians, five physicians, two genetic counsellors and two ethicists to identify a rare mutation in the XIAP that was causing widespread problems.
Due to recent cost reductions (see above) whole genome sequencing has become a realistic application in DNA diagnostics. In 2013, the 3Gb-TEST consortium obtained funding from the European Union to prepare the health care system for these innovations in DNA diagnostics. Quality assessment schemes, Health technology assessment and guidelines have to be in place. The 3Gb-TEST consortium has identified the analysis and interpretation of sequence data as the most complicated step in the diagnostic process. At the Consortium meeting in Athens in September 2014, the Consortium coined the word genotranslation for this crucial step. This step leads to a so-called genoreport. Guidelines are needed to determine the required content of these reports.
Genomes2People (G2P), an initiative of Brigham and Women's Hospital and Harvard Medical School was created in 2011 to examine the integration of genomic sequencing into clinical care of adults and children. G2P's director, Robert C. Green, had previously led the REVEAL study — Risk EValuation and Education for Alzheimer's Disease – a series of clinical trials exploring patient reactions to the knowledge of their genetic risk for Alzheimer's.
In 2018, researchers at Rady Children's Hospital Institute for Genomic Medicine in San Diego determined that rapid whole-genome sequencing (rWGS) could diagnose genetic disorders in time to change acute medical or surgical management (clinical utility) and improve outcomes in acutely ill infants. In a retrospective cohort study of acutely ill inpatient infants in a regional children's hospital from July 2016-March 2017, forty-two families received rWGS for etiologic diagnosis of genetic disorders. The diagnostic sensitivity of rWGS was 43% (eighteen of 42 infants) and 10% (four of 42 infants) for standard genetic tests (P = .0005). The rate of clinical utility of rWGS (31%, thirteen of 42 infants) was significantly greater than for standard genetic tests (2%, one of 42; P = .0015). Eleven (26%) infants with diagnostic rWGS avoided morbidity, one had a 43% reduction in likelihood of mortality, and one started palliative care. In six of the eleven infants, the changes in management reduced inpatient cost by $800,000-$2,000,000. The findings replicated a prior study of the clinical utility of rWGS in acutely ill inpatient infants, and demonstrated improved outcomes, net healthcare savings and consideration as a first tier test in this setting.
A 2018 review of 36 publications found the cost for whole genome sequencing to range from $1,906USD to $24,810USD and have a wide variance in diagnostic yield from 17% to 73% depending on patient groups.
Rare variant association study
Whole genome sequencing studies enable the assessment of associations between complex traits and both coding and noncoding rare variants (minor allele frequency (MAF) < 1%) across the genome. Single-variant analyses typically have low power to identify associations with rare variants, and variant set tests have been proposed to jointly test the effects of given sets of multiple rare variants. SNP annotations help to prioritize rare functional variants, and incorporating these annotations can effectively boost the power of genetic association of rare variants analysis of whole genome sequencing studies. Some tools have been specifically developed to provide all-in-one rare variant association analysis for whole-genome sequencing data, including integration of genotype data and their functional annotations, association analysis, result summary and visualization.
Meta-analysis of whole genome sequencing studies provides an attractive solution to the problem of collecting large sample sizes for discovering rare variants associated with complex phenotypes. Some methods have been developed to enable functionally informed rare variant association analysis in biobank-scale cohorts using efficient approaches for summary statistic storage.
Oncology
In this field, whole genome sequencing represents a great set of improvements and challenges to be faced by the scientific community, as it makes it possible to analyze, quantify and characterize circulating tumor DNA (ctDNA) in the bloodstream. This serves as a basis for early cancer diagnosis, treatment selection and relapse monitoring, as well as for determining the mechanisms of resistance, metastasis and phylogenetic patterns in the evolution of cancer. It can also help in the selection of individualized treatments for patients suffering from this pathology and observe how existing drugs are working during the progression of treatment. Deep whole genome sequencing involves a subclonal reconstruction based on ctDNA in plasma that allows for complete epigenomic and genomic profiling, showing the expression of circulating tumor DNA in each case.
Newborn screening
In 2013, Green and a team of researchers launched the BabySeq Project to study the ethical and medical consequences of sequencing a newborn's DNA.
As of 2015, whole genome and exome sequencing as a newborn screening tool were deliberated and in 2021, further discussed.
In 2021, the NIH funded BabySeq2, an implementation study that expanded the BabySeq project, enrolling 500 infants from diverse families and track the effects of their genomic sequencing on their pediatric care.
In 2023, the Lancet opined that in the UK "focusing on improving screening by upgrading targeted gene panels might be more sensible in the short term. Whole genome sequencing in the long term deserves thorough examination and universal caution."
Ethical concerns
The introduction of whole genome sequencing may have ethical implications. On one hand, genetic testing can potentially diagnose preventable diseases, both in the individual undergoing genetic testing and in their relatives. On the other hand, genetic testing has potential downsides such as genetic discrimination, loss of anonymity, and psychological impacts such as discovery of non-paternity.
Some ethicists insist that the privacy of individuals undergoing genetic testing must be protected, and is of particular concern when minors undergo genetic testing. Illumina's CEO, Jay Flatley, wrongly claimed in February 2009 that "by 2019 it will have become routine to map infants' genes when they are born". This potential use of genome sequencing is highly controversial, as it runs counter to established ethical norms for predictive genetic testing of asymptomatic minors that have been well established in the fields of medical genetics and genetic counseling. The traditional guidelines for genetic testing have been developed over the course of several decades since it first became possible to test for genetic markers associated with disease, prior to the advent of cost-effective, comprehensive genetic screening.
When an individual undergoes whole genome sequencing, they reveal information about not only their own DNA sequences, but also about probable DNA sequences of their close genetic relatives. This information can further reveal useful predictive information about relatives' present and future health risks. Hence, there are important questions about what obligations, if any, are owed to the family members of the individuals who are undergoing genetic testing. In Western/European society, tested individuals are usually encouraged to share important information on any genetic diagnoses with their close relatives, since the importance of the genetic diagnosis for offspring and other close relatives is usually one of the reasons for seeking a genetic testing in the first place. Nevertheless, a major ethical dilemma can develop when the patients refuse to share information on a diagnosis that is made for serious genetic disorder that is highly preventable and where there is a high risk to relatives carrying the same disease mutation. Under such circumstances, the clinician may suspect that the relatives would rather know of the diagnosis and hence the clinician can face a conflict of interest with respect to patient-doctor confidentiality.
Privacy concerns can also arise when whole genome sequencing is used in scientific research studies. Researchers often need to put information on patient's genotypes and phenotypes into public scientific databases, such as locus specific databases. Although only anonymous patient data are submitted to locus specific databases, patients might still be identifiable by their relatives in the case of finding a rare disease or a rare missense mutation. Public discussion around the introduction of advanced forensic techniques (such as advanced familial searching using public DNA ancestry websites and DNA phenotyping approaches) has been limited, disjointed, and unfocused. As forensic genetics and medical genetics converge toward genome sequencing, issues surrounding genetic data become increasingly connected, and additional legal protections may need to be established.
Public human genome sequences
First people with public genome sequences
The first nearly complete human genomes sequenced were two Americans of predominantly Northwestern European ancestry in 2007 (J. Craig Venter at 7.5-fold coverage, and James Watson at 7.4-fold). This was followed in 2008 by sequencing of an anonymous Han Chinese man (at 36-fold), a Yoruban man from Nigeria (at 30-fold), a female clinical geneticist (Marjolein Kriek) from the Netherlands (at 7 to 8-fold), and a female leukemia patient in her mid-50s (at 33 and 14-fold coverage for tumor and normal tissues). Steve Jobs was among the first 20 people to have their whole genome sequenced, reportedly for the cost of $100,000. , there were 69 nearly complete human genomes publicly available. In November 2013, a Spanish family made their personal genomics data publicly available under a Creative Commons public domain license. The work was led by Manuel Corpas and the data obtained by direct-to-consumer genetic testing with 23andMe and the Beijing Genomics Institute. This is believed to be the first such Public Genomics dataset for a whole family.
Databases
According to Science, the major databases of whole genomes are:
Genomic coverage
In terms of genomic coverage and accuracy, whole genome sequencing can broadly be classified into either of the following:
A draft sequence, covering approximately 90% of the genome at approximately 99.9% accuracy
A finished sequence, covering more than 95% of the genome at approximately 99.99% accuracy
Producing a truly high-quality finished sequence by this definition is very expensive. Thus, most human "whole genome sequencing" results are draft sequences (sometimes above and sometimes below the accuracy defined above).
| Technology | Biotechnology | null |
3956618 | https://en.wikipedia.org/wiki/Join%20and%20meet | Join and meet | In mathematics, specifically order theory, the join of a subset of a partially ordered set is the supremum (least upper bound) of denoted and similarly, the meet of is the infimum (greatest lower bound), denoted In general, the join and meet of a subset of a partially ordered set need not exist. Join and meet are dual to one another with respect to order inversion.
A partially ordered set in which all pairs have a join is a join-semilattice. Dually, a partially ordered set in which all pairs have a meet is a meet-semilattice. A partially ordered set that is both a join-semilattice and a meet-semilattice is a lattice. A lattice in which every subset, not just every pair, possesses a meet and a join is a complete lattice. It is also possible to define a partial lattice, in which not all pairs have a meet or join but the operations (when defined) satisfy certain axioms.
The join/meet of a subset of a totally ordered set is simply the maximal/minimal element of that subset, if such an element exists.
If a subset of a partially ordered set is also an (upward) directed set, then its join (if it exists) is called a directed join or directed supremum. Dually, if is a downward directed set, then its meet (if it exists) is a directed meet or directed infimum.
Definitions
Partial order approach
Let be a set with a partial order and let An element of is called the (or or ) of and is denoted by if the following two conditions are satisfied:
(that is, is a lower bound of ).
For any if then (that is, is greater than or equal to any other lower bound of ).
The meet need not exist, either since the pair has no lower bound at all, or since none of the lower bounds is greater than all the others. However, if there is a meet of then it is unique, since if both are greatest lower bounds of then and thus If not all pairs of elements from have a meet, then the meet can still be seen as a partial binary operation on
If the meet does exist then it is denoted If all pairs of elements from have a meet, then the meet is a binary operation on and it is easy to see that this operation fulfills the following three conditions: For any elements
(commutativity),
(associativity), and
(idempotency).
Joins are defined dually with the join of if it exists, denoted by
An element of is the (or or ) of in if the following two conditions are satisfied:
(that is, is an upper bound of ).
For any if then (that is, is less than or equal to any other upper bound of ).
Universal algebra approach
By definition, a binary operation on a set is a if it satisfies the three conditions a, b, and c. The pair is then a meet-semilattice. Moreover, we then may define a binary relation on A, by stating that if and only if In fact, this relation is a partial order on Indeed, for any elements
since by c;
if then by a; and
if then since then by b.
Both meets and joins equally satisfy this definition: a couple of associated meet and join operations yield partial orders which are the reverse of each other. When choosing one of these orders as the main ones, one also fixes which operation is considered a meet (the one giving the same order) and which is considered a join (the other one).
Equivalence of approaches
If is a partially ordered set, such that each pair of elements in has a meet, then indeed if and only if since in the latter case indeed is a lower bound of and since is the lower bound if and only if it is a lower bound. Thus, the partial order defined by the meet in the universal algebra approach coincides with the original partial order.
Conversely, if is a meet-semilattice, and the partial order is defined as in the universal algebra approach, and for some elements then is the greatest lower bound of with respect to since
and therefore
Similarly, and if is another lower bound of then whence
Thus, there is a meet defined by the partial order defined by the original meet, and the two meets coincide.
In other words, the two approaches yield essentially equivalent concepts, a set equipped with both a binary relation and a binary operation, such that each one of these structures determines the other, and fulfill the conditions for partial orders or meets, respectively.
Meets of general subsets
If is a meet-semilattice, then the meet may be extended to a well-defined meet of any non-empty finite set, by the technique described in iterated binary operations. Alternatively, if the meet defines or is defined by a partial order, some subsets of indeed have infima with respect to this, and it is reasonable to consider such an infimum as the meet of the subset. For non-empty finite subsets, the two approaches yield the same result, and so either may be taken as a definition of meet. In the case where subset of has a meet, in fact is a complete lattice; for details, see completeness (order theory).
Examples
If some power set is partially ordered in the usual way (by ) then joins are unions and meets are intersections; in symbols, (where the similarity of these symbols may be used as a mnemonic for remembering that denotes the join/supremum and denotes the meet/infimum).
More generally, suppose that is a family of subsets of some set that is partially ordered by
If is closed under arbitrary unions and arbitrary intersections and if belong to then
But if is not closed under unions then exists in if and only if there exists a unique -smallest such that
For example, if then whereas if then does not exist because the sets are the only upper bounds of in that could possibly be the upper bound but and
If then does not exist because there is no upper bound of in
| Mathematics | Order theory | null |
3957510 | https://en.wikipedia.org/wiki/Capacity%20factor | Capacity factor | The net capacity factor is the unitless ratio of actual electrical energy output over a given period of time to the theoretical maximum electrical energy output over that period. The theoretical maximum energy output of a given installation is defined as that due to its continuous operation at full nameplate capacity over the relevant period. The capacity factor can be calculated for any electricity producing installation, such as a fuel consuming power plant or one using renewable energy, such as wind, the sun or hydro-electric installations. The average capacity factor can also be defined for any class of such installations, and can be used to compare different types of electricity production.
The actual energy output during that period and the capacity factor vary greatly depending on a range of factors. The capacity factor can never exceed the availability factor, or uptime during the period. Uptime can be reduced due to, for example, reliability issues and maintenance, scheduled or unscheduled. Other factors include the design of the installation, its location, the type of electricity production and with it either the fuel being used or, for renewable energy, the local weather conditions. Additionally, the capacity factor can be subject to regulatory constraints and market forces, potentially affecting both its fuel purchase and its electricity sale.
The capacity factor is often computed over a timescale of a year, averaging out most temporal fluctuations. However, it can also be computed for a month to gain insight into seasonal fluctuations. Alternatively, it can be computed over the lifetime of the power source, both while operational and after decommissioning. A capacity factor can also be expressed and converted to full load hours.
Definition
Sample calculations
Nuclear power plant
Nuclear power plants are at the high end of the range of capacity factors, ideally reduced only by the availability factor, i.e. maintenance and refueling. The largest nuclear plant in the US, Palo Verde Nuclear Generating Station has between its three reactors a nameplate capacity of 3,942 MW. In 2010 its annual generation was 31,200,000 MWh, leading to a capacity factor of:
Each of Palo Verde’s three reactors is refueled every 18 months, with one refueling every spring and fall. In 2014, a refueling was completed in a record 28 days, compared to the 35 days of downtime that the 2010 capacity factor corresponds to.
In 2019, Prairie Island 1 was the US unit with the highest factor and actually reached 104.4%.
Wind farm
The Danish offshore wind farm Horns Rev 2 has a nameplate capacity of 209.3 MW.
it has produced 6416 GWh since its commissioning 7 years ago, i.e. an average annual production of 875 GWh/year and a capacity factor of:
Sites with lower capacity factors may be deemed feasible for wind farms, for example the onshore 1 GW Fosen Vind which is under construction in Norway has a projected capacity factor of 39%. Feasibility calculations may be affected by seasonality. For example in Finland, capacity factor during the cold winter months is more than double compared to July. While the annual average in Finland is 29.5%, the high demand for heating energy correlates with the higher capacity factor during the winter.
Certain onshore wind farms can reach capacity factors of over 60%, for example the 44 MW Eolo plant in Nicaragua had a net generation of 232.132 GWh in 2015, equivalent to a capacity factor of 60.2%, while United States annual capacity factors from 2013 through 2016 range from 32.2% to 34.7%.
Since the capacity factor of a wind turbine measures actual production relative to possible production, it is unrelated to Betz's coefficient of 16/27 59.3%, which limits production vs. energy available in the wind.
Hydroelectric dam
Three Gorges Dam in China is, with its nameplate capacity of 22,500 MW, the largest power generating station in the world by installed capacity.
In 2015 it generated 87 TWh, for a capacity factor of:
Hoover Dam has a nameplate capacity of 2080 MW and an annual generation averaging 4.2 TW·h. (The annual generation has varied between a high of 10.348 TW·h in 1984, and a low of 2.648 TW·h in 1956.).
Taking the average figure for annual generation gives a capacity factor of:
Photovoltaic power station
At the low range of capacity factors is the photovoltaic power station, which supplies power to the electricity grid from a large-scale photovoltaic system (PV system).
An inherent limit to its capacity factor comes from its requirement of daylight, preferably with a sun unobstructed by clouds, smoke or smog, shade from trees and building structures.
Since the amount of sunlight varies both with the time of the day and the seasons of the year, the capacity factor is typically computed on an annual basis.
The amount of available sunlight is mostly determined by the latitude of the installation and the local cloud cover.
The actual production is also influenced by local factors such as dust and ambient temperature, which ideally should be low. As for any power station, the maximum possible power production is the nameplate capacity times the number of hours in a year, while the actual production is the amount of electricity delivered annually to the grid.
For example, Agua Caliente Solar Project, located in Arizona near the 33rd parallel and awarded for its excellence in renewable energy has a nameplate capacity of 290 MW and an actual average annual production of 740 GWh/year.
Its capacity factor is thus:
.
A significantly lower capacity factor is achieved by Lauingen Energy Park located in Bavaria, near the 49th parallel. With a nameplate capacity of 25.7 MW and an actual average annual production of 26.98 GWh/year it has a capacity factor of 12.0%.
Determinants of a plant capacity factor
There are several reasons why a plant would have a capacity factor lower than 100%. These include technical constraints, such as operational availability of the plant (known as availability factor for electricity generators), economic reasons, and availability of an energy resource.
A plant can be out of service or operating at reduced output for part of the time due to equipment failures or routine maintenance.
This accounts for most of the unused capacity of base load power plants. Base load plants usually have low costs per unit of electricity because they are designed for maximum efficiency and are operated continuously at high output.
Geothermal power plants, nuclear power plants, coal-fired plants and bioenergy plants that burn solid material are almost always operated as base load plants, as they can be difficult to adjust to suit demand.
A plant can also have its output curtailed or intentionally left idle because the electricity is not needed or because the price of electricity is too low to make production economical.
This accounts for most of the unused capacity of peaking power plants and load following power plants.
Peaking plants may operate for only a few hours per year or up to several hours per day.
Many other power plants operate only at certain times of the day or year because of variation in loads and electricity prices.
If a plant is only needed during the day, for example, even if it operates at full power output from 8 am to 8 pm every day (12 hours) all year long, it would only have a 50% capacity factor.
Due to low capacity factors, electricity from peaking power plants is relatively expensive because the limited generation has to cover the plant fixed costs.
A third reason is that a plant may not have the fuel available to operate all of the time.
This can apply to fossil generating stations with restricted fuels supplies, but most notably applies to intermittent renewable resources.
Solar PV and wind turbines have a capacity factor limited by the availability of their "fuel", sunshine and wind respectively.
A hydroelectricity plant may have a capacity factor lower than 100% due to restriction or scarcity of water, or its output may be regulated to match the current power need, conserving its stored water for later usage.
Governments differ in their wiliness to accept risks of power outages and lack of resilience against natural disasters and military attack on electricity grids. Examples of historical events impacting grid resilience are the 1991 Gulf War air campaign against civilian infrastructure, 2015 Ukraine power grid hack, 2021 Texas power crisis and Russian strikes against Ukrainian infrastructure (2022–present). Low risk tolerance may require electricity grids to be more significantly overbuilt to mitigate the potential costs of electricity grid interruptions and outages, impacting on a technology-by-technology basis the amount of generation curtailment necessary under normal grid conditions.
Other reasons that a power plant may not have a capacity factor of 100% include restrictions or limitations on air permits and limitations on transmission that force the plant to curtail output.
Capacity factor of renewable energy
For renewable energy sources such as solar power, wind power and hydroelectricity, the main reason for reduced capacity factor is generally the availability of the energy source.
The plant may be capable of producing electricity, but its "fuel" (wind, sunlight or water) may not be available.
A hydroelectric plant's production may also be affected by requirements to keep the water level from getting too high or low and to provide water for fish downstream.
However, solar, wind and hydroelectric plants do have high availability factors, so when they have fuel available, they are almost always able to produce electricity.
When hydroelectric plants have water available, they are also useful for load following, because of their high dispatchability. A typical hydroelectric plant's operators can bring it from a stopped condition to full power in just a few minutes.
Wind farms are variable, due to the natural variability of the wind.
For a wind farm, the capacity factor is determined by the availability of wind, the swept area of the turbine and the size of the generator.
Transmission line capacity and electricity demand also affect the capacity factor.
Typical capacity factors of current wind farms are between 25 and 45%. In the United Kingdom during the five year period from 2011 to 2019 the annual capacity factor for wind was over 30%.
Solar energy is variable because of the daily rotation of the earth, seasonal changes, and because of cloud cover.
For example, the Sacramento Municipal Utility District observed a 15% capacity factor in 2005.
However, according to the SolarPACES programme of the International Energy Agency (IEA), solar power plants designed for solar-only generation are well matched to summer noon peak loads in areas with significant cooling demands, such as Spain or the south-western United States, although in some locations solar PV does not reduce the need for generation of network upgrades given that air conditioner peak demand often occurs in the late afternoon or early evening when solar output is reduced.
SolarPACES states that by using thermal energy storage systems the operating periods of solar thermal power (CSP) stations can be extended to become dispatchable (load following).
Geothermal has a higher capacity factor than many other power sources, and geothermal resources are generally available all the time.
Capacity factors by energy source
Worldwide
Nuclear power 88.7% (2006 - 2012 average of US's plants).
Hydroelectricity, worldwide average 44%, range of 10% - 99% depending on water availability (with or without regulation via storage dam).
Wind farms 21-52% (as of 2022).
CSP solar with storage and Natural Gas backup in Spain 63%, California 33%.
Photovoltaic solar in Germany 10%, Arizona 19%, Massachusetts 13.35% (8 year average as of July 2018).
United States
According to the US Energy Information Administration (EIA), from 2013 to 2017 the capacity factors of utility-scale generators were as follows:
|-
! colspan="8"| Non-fossil fuels !! Coal !! colspan="4" | Natural Gas !! colspan="3" | Petroleum Liquids
|-
!Nuclear !! Hydro !!Wind !!Solar PV !!Solar CSP !!Landfill Gas and !!Other Biomass including Wood !! Geothermal !! !! !! !! !! !! !! !!
|-
|89.9% ||38.9% ||32.4% ||NA ||NA ||68.9% ||56.7% ||73.6% ||59.8% ||48.2% ||4.9% ||10.6% ||6.1% ||12.1% ||0.8% ||2.2%
|-
|91.7% ||37.3% ||34.0% ||25.9% ||19.8% ||68.9% ||58.9% ||74.0% ||61.1% ||48.3% ||5.2% ||10.4% ||8.5% ||12.5% ||1.1% ||1.4%
|-
|92.3% ||35.8% ||32.2% ||25.8% ||22.1% ||68.7% ||55.3% ||74.3% ||54.7% ||55.9% ||6.9% ||11.5% ||8.9% ||13.3% ||1.1% ||2.2%
|-
|92.3% ||38.2% ||34.5% ||25.1% ||22.2% ||69.7% ||55.6% ||73.9% ||53.3% ||55.5% ||8.3% ||12.4% ||9.6% ||11.5% ||1.1% ||2.6%
|-
|92.2% ||43.1% ||34.6% ||25.7% ||21.8% ||68.0% ||57.8% ||74.0% ||53.7% ||51.3% ||6.7% ||10.5% || 9.9% ||13.5% ||0.9% || 2.3%
|-
|92.6% ||42.8% ||37.4% ||26.1% ||23.6% ||73.3% ||49.3% ||77.3% ||54.0% ||57.6% ||11.8% ||13.7% ||NA ||13.9% ||2.5% ||NA
However, these values often vary significantly by month.
United Kingdom
The following figures were collected by the Department of Energy and Climate Change on the capacity factors for various types of plants in UK grid:
| Technology | Concepts | null |
3957700 | https://en.wikipedia.org/wiki/Tritiated%20water | Tritiated water | Tritiated water is a radioactive form of water in which the usual protium atoms are replaced with tritium atoms. In its pure form it may be called tritium oxide (T2O or 3H2O) or super-heavy water. Pure T2O is a colorless liquid, and it is corrosive due to self-radiolysis. Diluted, tritiated water is mainly H2O plus some HTO (3HOH). It is also used as a tracer for water transport studies in life-science research. Furthermore, since it naturally occurs in minute quantities, it can be used to determine the age of various water-based liquids, such as vintage wines.
The name super-heavy water helps distinguish the tritiated material from heavy water, which contains deuterium instead.
Applications
Tritiated water can be used to measure an organism's total body water (TBW). Unlike doubly labeled water this method relies on scintillation counting. Tritiated water distributes itself into all body compartments relatively quickly. The concentration of tritiated water in urine is assumed to be similar to the concentration of tritiated water in the body. TBW is determined from the following relation:
Health risks
Tritium is radioactive and a low energy beta emitter.
While HTO is produced naturally by cosmic ray interactions in the stratosphere, it is also produced by human activities and can increase local concentrations and be considered an air and water pollutant. Anthropogenic sources of tritiated water include nuclear weapons testing, nuclear power plants, nuclear reprocessing and consumer products such as self-illuminating watches and signs.
HTO has a short biological half-life in the human body of 7 to 14 days, which both reduces the total effects of single-incident ingestion and precludes long-term bioaccumulation of HTO from the environment. The biological half life of tritiated water in the human body, which is a measure of body water turn-over, varies with the season. Studies on the biological half life of occupational radiation workers for free water tritium in a coastal region of Karnataka, India, show that the biological half life in the winter season is twice that of the summer season.
If tritium exposure is suspected or known, drinking uncontaminated water will help replace the tritium from the body. Increasing sweating, urination or breathing can help the body expel water and thereby the tritium contained in it. However, care should be taken that neither dehydration nor a depletion of the body's electrolytes results as the health consequences of those things (particularly in the short term) can be more severe than those of tritium exposure.
| Physical sciences | Water | Chemistry |
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