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2900052 | https://en.wikipedia.org/wiki/Calorimeter%20%28particle%20physics%29 | Calorimeter (particle physics) | In experimental particle physics, a calorimeter is a type of detector that measures the energy of particles. Particles enter the calorimeter and initiate a particle shower in which their energy is deposited in the calorimeter, collected, and measured. The energy may be measured in its entirety, requiring total containment of the particle shower, or it may be sampled. Typically, calorimeters are segmented transversely to provide information about the direction of the particle or particles, as well as the energy deposited, and longitudinal segmentation can provide information about the identity of the particle based on the shape of the shower as it develops. Calorimetry design is an active area of research in particle physics.
Types of calorimeters
Electromagnetic versus hadronic
such as electrons, positrons and photons. A . (See types of particle showers for the differences between the two.) Calorimeters are characterized by the radiation length (for ECALs) and nuclear interaction length (for HCALs) of their active material. ECALs tend to be 15–30 radiation lengths deep while HCALs are 5–8 nuclear interaction lengths deep.
Homogeneous versus sampling
An ECAL or an HCAL can be either a sampling calorimeter or a homogeneous calorimeter.
Calorimeters in high-energy physics experiments
Most particle physics experiments use some form of calorimetry. Often it is the most practical way to detect and measure neutral particles from an interaction. In addition, calorimeters are necessary for calculating "missing energy" which can be attributed to particles that rarely interact with matter and escape the detector, such as neutrinos. In most experiments the calorimeter works in conjunction with other components like a central tracker and a muon detector. All the detector components work together to achieve the objective of reconstructing a physics event.
| Physical sciences | Devices | Physics |
2902373 | https://en.wikipedia.org/wiki/Prehensile-tailed%20porcupine | Prehensile-tailed porcupine | The prehensile-tailed porcupines or coendous (genus Coendou) are found in Central and South America. Two other formerly recognized Neotropical tree porcupine genera, Echinoprocta and Sphiggurus, have been subsumed into Coendou, since Sphiggurus was shown by genetic studies to be polyphyletic, while Echinoprocta nested within Coendou.
Characteristics
Among the most notable features of Coendou porcupines are their unspined prehensile tails. The front and hind feet are also modified for grasping. These limbs all contribute to making this animal an adept climber, an adaptation to living most of their lives in trees.
They feed on leaves, shoots, fruits, bark, roots, and buds. They can be pests of plantation crops. They also make a distinctive "baby-like" sound to communicate in the wild.
Their young are born with soft hair that hardens to quills with age. Adults are slow-moving and will roll into a ball when threatened and on the ground. The record longevity is 27 years.
Species
Genus Coendou - prehensile-tailed porcupines
Baturite porcupine - C. baturitensis - a newly discovered species
Bicolored-spined porcupine - C. bicolor
Streaked dwarf porcupine - C. ichillus
Bahia porcupine - C. insidiosus
Black-tailed hairy dwarf porcupine - C. melanurus
Mexican hairy dwarf porcupine - C. mexicanus
Black dwarf porcupine - C. nycthemera
Amazonian long-tailed porcupine - C. longicaudatus
Brazilian porcupine - C. prehensilis
Frosted hairy dwarf porcupine - C. pruinosus
Andean porcupine - C. quichua
Roosmalen's dwarf porcupine - C. roosmalenorum
Stump-tailed porcupine - C. rufescens
Santa Marta porcupine - C. sanctamartae
C. speratus - a newly discovered species
Paraguaian hairy dwarf porcupine - C. spinosus
Brown hairy dwarf porcupine - C. vestitus
| Biology and health sciences | Rodents | Animals |
2903625 | https://en.wikipedia.org/wiki/Nadir | Nadir | The nadir is the direction pointing directly below a particular location; that is, it is one of two vertical directions at a specified location, orthogonal to a horizontal flat surface.
The direction opposite of the nadir is the zenith.
Definitions
Space science
Since the concept of being below is itself somewhat vague, scientists define the nadir in more rigorous terms. Specifically, in astronomy, geophysics and related sciences (e.g., meteorology), the nadir at a given point is the local vertical direction pointing in the direction of the force of gravity at that location.
The term can also be used to represent the lowest point that a celestial object reaches along its apparent daily path around a given point of observation (i.e. the object's lower culmination). This can be used to describe the position of the Sun, but it is only technically accurate for one latitude at a time and only possible at the low latitudes. The Sun is said to be at the nadir at a location when it is at the zenith at the location's antipode and is 90° below the horizon.
Nadir also refers to the downward-facing viewing geometry of an orbiting satellite, such as is employed during remote sensing of the atmosphere, as well as when an astronaut faces the Earth while performing a spacewalk. A nadir image is a satellite image or aerial photo of the Earth taken vertically. A satellite ground track represents its orbit projected to nadir on to Earth's surface.
Medicine
Generally in medicine, nadir is used to indicate the progression to the lowest point of a clinical symptom (e.g. fever patterns) or a laboratory count. In oncology, the term nadir is used to represent the lowest level of a blood cell count while a patient is undergoing chemotherapy. A diagnosis of neutropenic nadir after chemotherapy typically lasts 7–10 days.
Figurative usage
The word is also used figuratively to mean a low point, such as with a person's spirits, the quality of an activity or profession, or the nadir of American race relations.
| Physical sciences | Celestial sphere: General | Astronomy |
21541335 | https://en.wikipedia.org/wiki/Nitrososphaerota | Nitrososphaerota | The Nitrososphaerota (syn. Thaumarchaeota) are a phylum of the Archaea proposed in 2008 after the genome of Cenarchaeum symbiosum was sequenced and found to differ significantly from other members of the hyperthermophilic phylum Thermoproteota (formerly Crenarchaeota). Three described species in addition to C. symbiosum are Nitrosopumilus maritimus, Nitrososphaera viennensis, and Nitrososphaera gargensis. The phylum was proposed in 2008 based on phylogenetic data, such as the sequences of these organisms' ribosomal RNA genes, and the presence of a form of type I topoisomerase that was previously thought to be unique to the eukaryotes. This assignment was confirmed by further analysis published in 2010 that examined the genomes of the ammonia-oxidizing archaea Nitrosopumilus maritimus and Nitrososphaera gargensis, concluding that these species form a distinct lineage that includes Cenarchaeum symbiosum. The lipid crenarchaeol has been found only in Nitrososphaerota, making it a potential biomarker for the phylum. Most organisms of this lineage thus far identified are chemolithoautotrophic ammonia-oxidizers and may play important roles in biogeochemical cycles, such as the nitrogen cycle and the carbon cycle. Metagenomic sequencing indicates that they constitute ~1% of the sea surface metagenome across many sites.
Nitrososphaerota-derived membrane-spanning tetraether lipids (glycerol dialkyl glycerol tetraethers; GDGTs) from marine sediments can be used to reconstruct past temperatures via the TEX86 paleotemperature proxy, as these lipids vary in structure according to temperature. Because most Nitrososphaerota seem to be autotrophs that fix CO2, their GDGTs can act as a record for past Carbon-13 ratios in the dissolved inorganic carbon pool, and thus have the potential to be used for reconstructions of the carbon cycle in the past.
Taxonomy
The currently accepted taxonomy is based on the List of Prokaryotic names with Standing in Nomenclature (LPSN) and National Center for Biotechnology Information (NCBI)
Class Nitrososphaeria Stieglmeier et al. 2014 [Conexivisphaeria Kato et al. 2020]
?"Cenoporarchaeum" corrig. Zhang et al. 2019
?"Candidatus Gagatemarchaeum" Sheridan et al. 2023
?"Candidatus Giganthauma" Muller et al. 2010
?"Candidatus Nitrosodeserticola" Hwang et al. 2021
?"Candidatus Subgagatemarchaeum" Sheridan et al. 2023
Order "Geothermarchaeales" Adam et al. 2022
Family Geothermarchaeaceae Adam et al. 2022
"Geothermarchaeum" Adam et al. 2022
"Scotarchaeum" Adam et al. 2022
Order Conexivisphaerales Kato et al. 2020
Family Conexivisphaeraceae Kato et al. 2020
Conexivisphaera Kato et al. 2020
Order "Nitrosocaldales" de la Torre et al. 2008
Family "Nitrosocaldaceae" Qin et al. 2016
"Candidatus Nitrosothermus" Luo et al. 2021
"Candidatus Nitrosocaldus" de la Torre et al. 2008
Order Nitrososphaerales Stieglmeier et al. 2014
Family Methylarchaeaceae Hua et al. 2019
?"Candidatus Methylarchaeum" Hua et al. 2019
?"Candidatus Methanotowutia" Ou et al. 2022
Family Nitrososphaeraceae Stieglmeier et al. 2014
"Candidatus Nitrosocosmicus" Lehtovirta-Morley et al. 2016
"Candidatus Nitrosopolaris" Pessi, Rutanen & Hultman 2022
Nitrososphaera Stieglmeier et al. 2014
Order Nitrosopumilales Qin et al. 2017
Family Nitrosopumilaceae Qin et al. 2017
"Cenarchaeum" DeLong & Preston 1996
Nitrosarchaeum corrig. Jung et al. 2018
?"Candidatus Nitrosoabyssus" Garritano et al. 2024
?"Candidatus Nitrosokoinonia" Glasl et al. 2023
?"Candidatus Nitrosomaritimum" Zhao et al. 2024
"Candidatus Nitrosopelagicus" Santoro et al. 2015
Nitrosopumilus Qin et al. 2017
?"Candidatus Nitrosospongia" Moeller et al. 2019
"Candidatus Nitrosotalea" Lehtovirta 2011
"Candidatus Nitrosotenuis" Li et al. 2016
Metabolism
Nitrososphaerota are important ammonia oxidizers in aquatic and terrestrial environments, and are the first archaea identified as being involved in nitrification. They are capable of oxidizing ammonia at much lower substrate concentrations than ammonia-oxidizing bacteria, and so probably dominate in oligotrophic conditions. Their ammonia oxidation pathway requires less oxygen than that of ammonia-oxidizing bacteria, so they do better in environments with low oxygen concentrations like sediments and hot springs. Ammonia-oxidizing Nitrososphaerota can be identified metagenomically by the presence of archaeal ammonia monooxygenase (amoA) genes, which indicate that they are overall more dominant than ammonia oxidizing bacteria. In addition to ammonia, at least one Nitrososphaerota strain has been shown to be able to use urea as a substrate for nitrification. This would allow for competition with phytoplankton that also grow on urea. One study of microbes from wastewater treatment plants found that not all Nitrososphaerota that express amoA genes are active ammonia oxidizers. These Nitrososphaerota may be capable of oxidizing methane instead of ammonia, or they may be heterotrophic, indicating a potential for a diversity of metabolic lifestyles within the phylum. Marine Nitrososphaerota have also been shown to produce nitrous oxide, which as a greenhouse gas has implications for climate change. Isotopic analysis indicates that most nitrous oxide flux to the atmosphere from the ocean, which provides around 30% of the natural flux, may be due to the metabolic activities of archaea.
Many members of the phylum assimilate carbon by fixing HCO3−. This is done using a hydroxypropionate/hydroxybutyrate cycle similar to the Thermoproteota but which appears to have evolved independently. All Nitrososphaerota that have been identified by metagenomics thus far encode this pathway. Notably, the Nitrososphaerota CO2-fixation pathway is more efficient than any known aerobic autotrophic pathway. This efficiency helps explain their ability to thrive in low-nutrient environments. Some Nitrososphaerota such as Nitrosopumilus maritimus are able to incorporate organic carbon as well as inorganic, indicating a capacity for mixotrophy. At least two isolated strains have been identified as obligate mixotrophs, meaning they require a source of organic carbon in order to grow.
A study has revealed that Nitrososphaerota are most likely the dominant producers of the critical vitamin B12. This finding has important implications for eukaryotic phytoplankton, many of which are auxotrophic and must acquire vitamin B12 from the environment; thus the Nitrososphaerota could play a role in algal blooms and by extension global levels of atmospheric carbon dioxide. Because of the importance of vitamin B12 in biological processes such as the citric acid cycle and DNA synthesis, production of it by the Nitrososphaerota may be important for a large number of aquatic organisms.
Environment
Many Nitrososphaerota, such as Nitrosopumilus maritimus, are marine and live in the open ocean. Most of these planktonic Nitrososphaerota, which compose the Marine Group I.1a, are distributed in the subphotic zone, between 100m and 350m. Other marine Nitrososphaerota live in shallower waters. One study has identified two novel Nitrososphaerota species living in the sulfidic environment of a tropical mangrove swamp. Of these two species, Candidatus Giganthauma insulaporcus and Candidatus Giganthauma karukerense, the latter is associated with Gammaproteobacteria with which it may have a symbiotic relationship, though the nature of this relationship is unknown. The two species are very large, forming filaments larger than ever before observed in archaea. As with many Nitrososphaerota, they are mesophilic. Genetic analysis and the observation that the most basal identified Nitrososphaerota genomes are from hot environments suggests that the ancestor of Nitrososphaerota was thermophilic, and mesophily evolved later.
| Biology and health sciences | Archaea | Plants |
2088910 | https://en.wikipedia.org/wiki/Maritime%20passenger%20terminal | Maritime passenger terminal | A passenger terminal is a structure in a port which services passengers boarding and leaving water vessels such as ferries, cruise ships and ocean liners. Depending on the types of vessels serviced by the terminal, it may be named (for example) ferry terminal, cruise terminal, marine terminal or maritime passenger terminal. As well as passengers, a passenger terminal sometimes has facilities for automobiles and other land vehicles to be picked up and dropped off by the water vessel.
Facilities
Passenger terminals may vary greatly in size. A small ferry terminal servicing a commuter ferry may just have the means to tie up the vessel and a waiting area for passengers. Even for a large, vehicle-carrying cross-sea ferry, the terminal at a small island location may be similar sized, with just a short ramp to enable vehicles to be driven onto the ferry.
Passengers may be loaded onto a ship from the wharf by a gangway or by a linkspan. Goods packed in containers may be driven onto the vessel by a vehicle which then detaches itself from the container and returns to shore.
If the passenger terminal handles vehicles (which is common especially in cross-sea ferry terminals), it will usually have the facilities, such as appropriate markings on the ground, to enable the vehicles to line up in an orderly manner. Vehicles may be driven off the ship directly, if the vessel is a Roll-on/roll-off ship.
Passenger terminals in large ports usually have passenger facilities comparable with medium-sized airports, including waiting areas, ticketing desks, luggage deposit and retrieval areas, and food, beverage and other retail outlets. Ferry terminals for international ferries, such as those crossing between the United Kingdom and continental Europe, also have customs and immigration inspection facilities and security control areas similar to an international airport.
Historically, the largest passenger terminals were located in major coastal cities servicing large ocean liners. With the demise of most ocean liners in the later half of the 20th century and the rise of cruise ship tourism in its stead, the largest passenger terminals today are those in "cruise home ports". In addition to extensive facilities to service passengers, these terminals must also be capable of handling the large amount of supplies required by large cruise ships and ocean liners.
Major passenger ports (such as the Port of Southampton) tend to have numerous docks and wharves, some with multiple berths, in order to handle more than one ship simultaneously. Some ports → a single, large passenger terminal to service multiple docks, while others have multiple terminal buildings, each servicing a dock or wharf, so that passengers can board vessels directly from the terminal.
Major passenger terminals
Australia
Overseas Passenger Terminal, in Sydney
Station Pier, in Melbourne
Brisbane International Cruise Terminal In Brisbane
Canada
Canada Place, in Vancouver
China
International Cruise Terminal, in Shanghai
Hong Kong
Kai Tak Cruise Terminal, in Kai Tak
Ocean Terminal, in Tsim Sha Tsui
Singapore
Marina Bay Cruise Centre
Taiwan
Kaohsiung Port Cruise Terminal, in Kaohsiung
United Kingdom
Queen Elizabeth II Terminal, Mayflower Terminal, City Terminal, Horizon Terminal and Ocean Terminal, in Southampton
United States
Brooklyn Cruise Terminal, in New York
Manhattan Cruise Terminal, in New York
Port Everglades cruise terminal, in Fort Lauderdale
| Technology | Coastal infrastructure | null |
2090039 | https://en.wikipedia.org/wiki/Geology%20of%20Mercury | Geology of Mercury | The geology of Mercury is the scientific study of the surface, crust, and interior of the planet Mercury. It emphasizes the composition, structure, history, and physical processes that shape the planet. It is analogous to the field of terrestrial geology. In planetary science, the term geology is used in its broadest sense to mean the study of the solid parts of planets and moons. The term incorporates aspects of geophysics, geochemistry, mineralogy, geodesy, and cartography.
Historically, Mercury has been the least understood of all the terrestrial planets in the Solar System. This stems largely from its proximity to the Sun which makes reaching it with spacecraft technically challenging and Earth-based observations difficult. For decades, the principal source of geologic information about Mercury came from the 2,700 images taken by the Mariner 10 spacecraft during three flybys of the planet from 1974 to 1975. These images covered about 45% of the planet’s surface, but many of them were unsuitable for detailed geologic investigation because of high sun angles which made it hard to determine surface morphology and topography. This dearth of information was greatly alleviated by the MErcury Surface, Space ENvironment, GEochemistry, and Ranging (MESSENGER) spacecraft which between 2008 and 2015 collected over 291,000 images covering the entire planet, along with a wealth of other scientific data. The European Space Agency's (ESA's) BepiColombo spacecraft, scheduled to go into orbit around Mercury in 2026, is expected to help answer many of the remaining questions about Mercury's geology.
Mercury's surface is dominated by impact craters, basaltic rock and smooth plains, many of them a result of flood volcanism, similar in some respects to the lunar maria, and locally by pyroclastic deposits. Other notable features include vents which appear to be the source of magma-carved valleys, often-grouped irregular-shaped depressions termed "hollows" that are believed to be the result of collapsed magma chambers, scarps indicative of thrust faulting, and mineral deposits (possibly ice) inside craters at the poles. Long thought to be geologically inactive, new evidence suggests there may still be some level of activity.
Mercury's density implies a solid iron-rich core that accounts for about 60% of its volume (75% of its radius). Mercury's magnetic equator is shifted nearly 20% of the planet's radius towards the north, the largest ratio of all planets. This shift suggests there being one or more iron-rich molten layers surrounding the core producing a dynamo effect similar to that of Earth. Additionally, the offset magnetic dipole may result in uneven surface weathering by the solar wind, knocking more surface particles up into the southern exosphere and transporting them for deposit in the north. Scientists are gathering telemetry to determine if such is the case.
Difficulties in exploration
Reaching Mercury from Earth poses significant technical challenges, because the planet orbits so much closer to the Sun than does the Earth. A Mercury-bound spacecraft launched from Earth must travel 91 million kilometers into the Sun's gravitational potential well. Starting from the Earth's orbital speed of 30 km/s, the change in velocity (delta-v) the spacecraft must make to enter into a Hohmann transfer orbit that passes near Mercury is large compared to other planetary missions. The potential energy liberated by moving down the Sun's potential well becomes kinetic energy; requiring another large delta-v to do anything other than rapidly pass by Mercury. In order to land safely or enter a stable orbit the spacecraft must rely entirely on rocket motors because Mercury has negligible atmosphere. A direct trip to Mercury actually requires more rocket fuel than that required to escape the Solar System completely. As a result, only three space probes, Mariner 10, MESSENGER, and BepiColombo have visited Mercury so far.
Furthermore, the space environment near Mercury is demanding, posing the double dangers to spacecraft of intense solar radiation and high temperatures.
Historically, a second obstacle has been that Mercury's period of rotation is a slow 58 Earth days, so that spacecraft flybys are restricted to viewing only a single illuminated hemisphere. Unfortunately, even though Mariner 10 space probe flew past Mercury three times during 1974 and 1975, it observed the same area during each pass. This was because Mariner 10's orbital period was almost exactly 3 sidereal Mercury days, and the same face of the planet was lit at each of the close approaches. As a result, less than 45% of the planet's surface was mapped.
Earth-based observations are made difficult by Mercury's constant proximity to the Sun. This has several consequences:
Whenever the sky is dark enough for viewing through telescopes, Mercury is always already near the horizon, where viewing conditions are poor anyway due to atmospheric factors.
The Hubble Space Telescope and other space observatories are usually prevented from pointing close to the Sun for safety reasons (Erroneously pointing such sensitive instruments at the Sun is likely to cause permanent damage).
Geological history
Like the Earth, Moon and Mars, Mercury's geologic history is divided into eras. From oldest to youngest, these are: the pre-Tolstojan, Tolstojan, Calorian, Mansurian, and Kuiperian. Their ages are based on relative dating only.
After the formation of Mercury along with the rest of the Solar System 4.6 billion years ago, heavy bombardment by asteroids and comets ensued. The last intense bombardment phase, the Late Heavy Bombardment came to an end about 3.8 billion years ago. Some regions or massifs, a prominent one being the one that formed the Caloris Basin, were filled by magma eruptions from within the planet. These created smooth intercrater plains similar to the maria found on the Moon.
Later, as the planet cooled and contracted, its surface began to crack and form ridges; these surface cracks and ridges can be seen on top of other features, such as the craters and smoother plains—a clear indication that they are more recent.
Mercury's period of volcanism ended when the planet's mantle had contracted enough to prevent further lava from breaking through to the surface. This probably occurred at some point during its first 700 or 800 million years of history.
Since then, the main surface processes have been intermittent impacts.
Timeline
Time unit: millions of years
Surface features
Mercury's surface is overall similar in appearance to that of the Moon, with extensive mare-like plains and heavily cratered terrains similar to the lunar highlands and made locally by accumulations of pyroclastic deposits.
Boulders
Mercury compared to the Moon has a rarity of boulders; about thirty times fewer boulders are found on Mercury than on the Moon. An explanation for this rarity is that the lifespan of Mercury's boulders can be less than the life span of boulders on the Moon (around 100 million years).
Boulders that have been found on Mercury as associated with fresh impact craters that are hundreds of meters in diameter or larger.
Impact basins and craters
Craters on Mercury range in diameter from small bowl-shaped craters to multi-ringed impact basins hundreds of kilometers across. They appear in all states of degradation, from relatively fresh rayed-craters, to highly degraded crater remnants. Mercurian craters differ subtly from Lunar craters – the extent of their ejecta blankets is much smaller, which is a consequence of the 2.5 times stronger surface gravity on Mercury.
The largest known crater is the enormous Caloris Basin, with a diameter of 1,550 km. A basin of comparable size, tentatively named Skinakas Basin had been postulated from low resolution Earth-based observations of the non-Mariner-imaged hemisphere, but has not been observed in MESSENGER imagery of the corresponding terrain. The impact which created the Caloris Basin was so powerful that its effects are seen on a global scale. It caused lava eruptions and left a concentric ring over 2 km tall surrounding the impact crater. At the antipode of the Caloris Basin lies a large region of unusual, hilly and furrowed terrain, sometimes called "Weird Terrain". The favoured hypothesis for the origin of this geomorphologic unit is that shock waves generated during the impact traveled around the planet, and when they converged at the basin's antipode (180 degrees away) the high stresses were capable of fracturing the surface. A much less favoured idea was that this terrain formed as a result of the convergence of ejecta at this basin's antipode. Furthermore, the formation of the Caloris Basin appears to have produced a shallow depression concentric around the basin, which was later filled by the smooth plains (see below).
Overall about 15 impact basins have been identified on the imaged part of Mercury. Other notable basins include the 400 km wide, multi-ring, Tolstoj Basin which has an ejecta blanket extending up to 500 km from its rim, and its floor has been filled by smooth plains materials. Beethoven Basin also has a similar-sized ejecta blanket and a 625 km diameter rim.
As on the Moon, fresh craters on Mercury show prominent bright ray systems. These are made by ejected debris, which tend to be brighter while they remain relatively fresh because of a lesser amount of space weathering than the surrounding older terrain.
Pit-floor craters
Some impact craters on Mercury have non-circular, irregularly shaped depressions or pits on their floors. Such craters have been termed pit-floor craters, and MESSENGER team members have suggested that such pits formed by the collapse of subsurface magma chambers. If this suggestion is correct, the pits are evidence of volcanic processes at work on Mercury. Pit craters are rimless, often irregularly shaped, and steep-sided, and they display no associated ejecta or lava flows but are typically distinctive in color. For example, the pits of Praxiteles have an orange hue. Thought to be evidence of shallow magmatic activity, pit craters may have formed when subsurface magma drained elsewhere and left a roof area unsupported, leading to collapse and the formation of the pit. Major craters exhibiting these features include Beckett, Gibran, Lermontov, Picasso, and Navoi, among others. It was suggested that these pits with associated brighter and redder deposits may be pyroclastic deposits caused by explosive volcanism.
Plains
There are two geologically distinct plains units on Mercury:
Intercrater plains are the oldest visible surface, predating the heavily cratered terrain. They are gently rolling or hilly and occur in the regions between larger craters. The intercrater plains appear to have obliterated many earlier craters, and show a general paucity of smaller craters below about 30 km in diameter. It is not clear whether they are of volcanic or impact origin. The intercrater plains are distributed roughly uniformly over the entire surface of the planet.
Smooth plains are widespread flat areas resembling the lunar maria, which fill depressions of various sizes. Notably, they fill a wide ring surrounding the Caloris Basin. An appreciable difference to the lunar maria is that the smooth plains of Mercury have the same albedo as the older intercrater plains. Despite a lack of unequivocally volcanic features, their localisation and lobate-shaped colour units strongly support a volcanic origin. All the Mercurian smooth plains formed significantly later than the Caloris basin, as evidenced by appreciably smaller crater densities than on the Caloris ejecta blanket.
The floor of the Caloris Basin is also filled by a geologically distinct flat plain, broken up by ridges and fractures in a roughly polygonal pattern. It is not clear whether they are volcanic lavas induced by the impact, or a large sheet of impact melt.
Tectonic features
One unusual feature of the planet's surface is the numerous compression folds which crisscross the plains. It is thought that as the planet's interior cooled, it contracted and its surface began to deform. The folds can be seen on top of other features, such as craters and smoother plains, indicating that they are more recent. Mercury's surface is also flexed by significant tidal bulges raised by the Sun—the Sun's tides on Mercury are about 17% stronger than the Moon's on Earth.
Faculae
Faculae on Mercury are bright areas often surrounding irregular depressions. They are generally interpreted to be pyroclastic in nature. The faculae on Mercury are all named using words in different languages meaning snake.
Terminology
Non-crater surface features are given the following names:
Albedo features – areas of markedly different reflectivity
Dorsa — ridges (see List of ridges on Mercury)
Montes — mountains (see List of mountains on Mercury)
Planitiae — plains (see List of plains on Mercury)
Rupes — scarps (see List of scarps on Mercury)
Valles — valleys (see List of valleys on Mercury)
High-albedo polar patches and possible presence of ice
The first radar observations of Mercury were carried out by the radiotelescopes at Arecibo (Puerto Rico) and Goldstone (California, United States), with assistance from the U.S. National Radio Astronomy Observatory Very Large Array (VLA) facility in New Mexico. The transmissions sent from the NASA Deep Space Network site at Goldstone were at a power level of 460 kW at 8.51 GHz; the signals received by the VLA multi-dish array detected points of radar reflectivity (radar luminosity) with depolarized waves from Mercury's north pole.
Radar maps of the surface of the planet were made using the Arecibo radiotelescope. The survey was conducted with 420 kW UHF band (2.4 GHz) radio waves which allowed for a 15 km resolution. This study not only confirmed the existence of the zones of high reflectivity and depolarization, but also found a number of new areas (bringing the total to 20) and was even able to survey the poles. It has been postulated that surface ice may be responsible for these high luminosity levels, as the silicate rocks that compose most of the surface of Mercury have exactly the opposite effect on luminosity.
In spite of its proximity to the Sun, Mercury may have surface ice, since temperatures near the poles are constantly below freezing point: On the polar plains, the temperature does not rise above −106 °C. Craters at Mercury's higher latitudes (discovered by radar surveys from Earth as well) may be deep enough to shield the ice from direct sunlight. Inside the craters, where there is no solar light, temperatures fall to −171 °C.
Despite sublimation into the vacuum of space, the temperature in the permanently shadowed region is so low that this sublimation is slow enough to potentially preserve deposited ice for billions of years.
At the South Pole, the location of a large zone of high reflectivity coincides with the location of the Chao Meng-Fu crater, and other small craters containing reflective areas have also been identified. At the North Pole, a number of craters smaller than Chao-Meng Fu have these reflective properties.
The strength of the radar reflections seen on Mercury are small compared to that which would occur with pure ice. This may be due to powder deposition that does not cover the surface of the crater completely or other causes, e.g. a thin overlying surface layer. However, the evidence for ice on Mercury is not definitive. The anomalous reflective properties could also be due to the existence of deposits of metallic sulfates or other materials with high reflectance.
Possible origin of ice
Mercury is not unique in having craters that stand in permanent shadow; at the south pole of Earth's Moon there is a large crater (Aitken) where some possible signs of the presence of ice have been seen (although their interpretation is disputed). It is thought by astronomers that ice on both Mercury and the Moon must have originated from external sources, mostly impacting comets. These are known to contain large amounts, or a majority, of ice. It is therefore conceivable for meteorite impacts to have deposited water in the permanently shadow craters, where it would remain unwarmed for possibly billions of years due to the lack of an atmosphere to efficiently conduct heat and stable orientation of Mercury's rotation axis.
Mercury's biological history
Habitability
There may be scientific support, based on studies reported in March 2020, for considering that parts of the planet Mercury may have been habitable, and perhaps that life forms, albeit likely primitive microorganisms, may have existed on the planet.
| Physical sciences | Solar System | Astronomy |
2091630 | https://en.wikipedia.org/wiki/Iron%28II%29%20sulfide | Iron(II) sulfide | Iron(II) sulfide or ferrous sulfide (Br.E. sulphide) is one of a family of chemical compounds and minerals with the approximate formula . Iron sulfides are often iron-deficient non-stoichiometric. All are black, water-insoluble solids.
Preparation and structure
FeS can be obtained by the heating of iron and sulfur:
Fe + S → FeS
FeS adopts the nickel arsenide structure, featuring octahedral Fe centers and trigonal prismatic sulfide sites.
Reactions
Iron sulfide reacts with hydrochloric acid, releasing hydrogen sulfide:
FeS + 2 HCl → FeCl2 + H2S
FeS + H2SO4 → FeSO4 + H2S
In moist air, iron sulfides oxidize to hydrated ferrous sulfate.
Biology and biogeochemistry
Iron sulfides occur widely in nature in the form of iron–sulfur proteins.
As organic matter decays under low-oxygen (or hypoxic) conditions such as in swamps or dead zones of lakes and oceans, sulfate-reducing bacteria reduce various sulfates present in the water, producing hydrogen sulfide. Some of the hydrogen sulfide will react with metal ions in the water or solid to produce iron or metal sulfides, which are not water-soluble. These metal sulfides, such as iron(II) sulfide, are often black or brown, leading to the color of sludge.
Pyrrhotite is a waste product of the Desulfovibrio bacteria, a sulfate reducing bacteria.
When eggs are cooked for a long time, the yolk's surface may turn green. This color change is due to iron(II) sulfide, which forms as iron from the yolk reacts with hydrogen sulfide released from the egg white by the heat. This reaction occurs more rapidly in older eggs as the whites are more alkaline.
The presence of ferrous sulfide as a visible black precipitate in the growth medium peptone iron agar can be used to distinguish between microorganisms that produce the cysteine metabolizing enzyme cysteine desulfhydrase and those that do not. Peptone iron agar contains the amino acid cysteine and a chemical indicator, ferric citrate. The degradation of cysteine releases hydrogen sulfide gas that reacts with the ferric citrate to produce ferrous sulfide.
| Physical sciences | Sulfide salts | Chemistry |
5248591 | https://en.wikipedia.org/wiki/Shower%20gel | Shower gel | Shower gel (also called body wash) is a specialized liquid product used for cleaning the body during showers. Not to be confused with liquid soaps, shower gels, in fact, do not contain saponified oil. Instead, it uses synthetic detergents derived from either petroleum or plant sources.
Body washes and shower gels have a lower pH value than the traditional soap, which is also known to feel less drying to the skin. In certain cases, sodium stearate is added to the chemical combination to create a solid version of the shower gel.
History
Shower gel is a derivative invention of the liquid soap, which first appeared in the 1800s. In 1865, William Shepphard patented the formula behind the liquid soap, but the product gained eventual popularity with the rise of Palmolive soap in 1898, by B.J. Johnson.
Modern chemistry later enabled the creation of the shower gel, which specialized in cleaning the entire body during baths and showers.
Properties
Shower gels are known to consist of the same basic ingredients as soap - water, betaines, and sodium laureth sulfate, or SLS. But the main difference between the two products lie in its surfactants - compounds known to lower the surface tension between substances, which helps in the emulsification and the washing away of oily dirt. The surfactants of shower gels do not come from saponification, that is by reacting a type of oil or fat with lye. Instead, it uses synthetic detergents for surfactants derived from either plant-based sources or petroleum. This gives the product a lower pH value than soap and might also feel less drying to the skin. Some people have likened the effect to feeling less squeaky clean, however.
Surfactants can make up as much as 50 percent of the shower gel content, with the remaining proportion being made up of a combination of water and ingredients to thicken, preserve, emulsify, add fragrance, and color. Multiple surfactants are often used to achieve desired product qualities. A primary surfactant can provide good foaming ability and cleaning effectiveness, while a secondary surfactant can add qualities of mildness to prevent irritation or over-drying of the skin. To prevent shower gel ingredients from separating, emulsifiers such as diethanolamine are added. Conditioning agents may also be added to moisturize the skin during and after product use. They are also available in different colours and scents. Ingredients, like scent in the form of essential oils or fragrance oils and colorant in the form of water-soluble dyes are common in shower gels.
Microbeads were commonly used in shower gels until recently. Microbeads are tiny spheres of plastic that were added to a variety of cosmetic products for their exfoliating qualities. They are too small to filter out of water systems and end up in waterways and oceans, potentially passing toxins to animal life and humans. Following the legislative actions of other countries, the United States passed the Microbead-Free Waters Act in 2015, which bans microbeads in the U.S. incrementally starting in 2017, with full implementation set for 2019. It has been banned from production and use in cosmetics in the U.S. since July 1, 2017, and in the UK since October 1, 2018.
Shower gels for men may contain the ingredient menthol, which gives a cooling and stimulating sensation on the skin, and some men's shower gels are also designed specifically for use on hair and body. Shower gels contain milder surfactant bases than shampoos, and some also contain gentle conditioning agents in the formula. This means that shower gels can also double as an effective and perfectly acceptable substitute to shampoo, even if they are not labelled as a hair and body wash. Washing hair with shower gel should give approximately the same result as using a moisturising shampoo.
Marketing
Like shampoo and bubble bath products, many are marketed directly towards children. These often feature scents intended to appeal to children, such as fruit scents, or cookies or cotton candy scents. Many bottles feature popular characters from children's television or movies. As with men's body wash, they often are specifically designed to be used also as a shampoo and conditioner. They also often contain gentle ingredients designed for young skin.
| Biology and health sciences | Hygiene products | Health |
33158412 | https://en.wikipedia.org/wiki/Scholarly%20peer%20review | Scholarly peer review | Scholarly peer review or academic peer review (also known as refereeing) is the process of having a draft version of a researcher's methods and findings reviewed (usually anonymously) by experts (or "peers") in the same field. Peer review is widely used for helping the academic publisher (that is, the editor-in-chief, the editorial board or the program committee) decide whether the work should be accepted, considered acceptable with revisions, or rejected for official publication in an academic journal, a monograph or in the proceedings of an academic conference. If the identities of authors are not revealed to each other, the procedure is called dual-anonymous peer review.
Academic peer review requires a community of experts in a given (and often narrowly defined) academic field, who are qualified and able to perform reasonably impartial review. Impartial review, especially of work in less narrowly defined or inter-disciplinary fields, may be difficult to accomplish, and the significance (good or bad) of an idea may never be widely appreciated among its contemporaries. Peer review is generally considered necessary to academic quality and is used in most major scholarly journals. However, peer review does not prevent publication of invalid research, and as experimentally controlled studies of this process are difficult to arrange, direct evidence that peer review improves the quality of published papers is scarce.
History
The first record of an editorial pre-publication peer-review is from 1665 by Henry Oldenburg, the founding editor of Philosophical Transactions of the Royal Society at the Royal Society of London.
The first peer-reviewed publication might have been the Medical Essays and Observations published by the Royal Society of Edinburgh in 1731. The present-day peer-review system evolved from this 18th-century process, began to involve external reviewers in the mid-19th-century, and did not become commonplace until the mid-20th-century.
Peer review became a touchstone of the scientific method, but until the end of the 19th century was often performed directly by an editor-in-chief or editorial committee. Editors of scientific journals at that time made publication decisions without seeking outside input, i.e. an external panel of reviewers, giving established authors latitude in their journalistic discretion. For example, Albert Einstein's four revolutionary Annus Mirabilis papers in the 1905 issue of Annalen der Physik were evaluated by the journal's editor-in-chief, Max Planck, and its co-editor, Wilhelm Wien, both future Nobel prize winners and together experts on the topics of these papers. On a much later occasion, Einstein was severely critical of the external review process, saying that he had not authorized the editor in chief to show his manuscript "to specialists before it is printed", and informing him that he would "publish the paper elsewhere"which he did, with substantial modifications.
While some medical journals started to systematically appoint external reviewers, it is only since the middle of the 20th century that this practice has spread widely and that external reviewers have been given some visibility within academic journals, including being thanked by authors and editors. A 2003 editorial in Nature stated that, in the early 20th century, "the burden of proof was generally on the opponents rather than the proponents of new ideas." Nature itself instituted formal peer review only in 1967. Journals such as Science and the American Journal of Medicine increasingly relied on external reviewers in the 1950s and 1960s, in part to reduce the editorial workload. In the 20th century, peer review also became common for science funding allocations. This process appears to have developed independently from that of editorial peer review.
Gaudet provides a social science view of the history of peer review carefully tending to what is under investigation, here peer review, and not only looking at superficial or self-evident commonalities among inquisition, censorship, and journal peer review. It builds on historical research by Gould, Biagioli, Spier, and Rip. The first Peer Review Congress met in 1989. Over time, the fraction of papers devoted to peer review has steadily declined, suggesting that as a field of sociological study, it has been replaced by more systematic studies of bias and errors. In parallel with "common experience" definitions based on the study of peer review as a "pre-constructed process", some social scientists have looked at peer review without considering it as pre-constructed. Hirschauer proposed that journal peer review can be understood as reciprocal accountability of judgements among peers. Gaudet proposed that journal peer review could be understood as a social form of boundary judgementdetermining what can be considered as scientific (or not) set against an overarching knowledge system, and following predecessor forms of inquisition and censorship.
Pragmatically, peer review refers to the work done during the screening of submitted manuscripts. This process encourages authors to meet the accepted standards of their discipline and reduces the dissemination of irrelevant findings, unwarranted claims, unacceptable interpretations, and personal views. Publications that have not undergone peer review are likely to be regarded with suspicion by academic scholars and professionals. Non-peer-reviewed work does not contribute, or contributes less, to the academic credit of a scholar (such as the h-index), although this heavily depends on the field.
Justification
It is difficult and time-consuming for authors and researchers, whether individually or in a team, to spot and provide feedback on every mistake or flaw in a complicated piece of work. This is not necessarily a reflection on those concerned, but because with a new and perhaps eclectic subject, an opportunity for improvement may be more obvious to someone with special expertise or who simply looks at it with a fresh eye. Therefore, showing work to others increases the probability that weaknesses will be identified and improved. For both grant-funding and publication in a scholarly journal, it is also normally a requirement that the subject is both novel and substantial.
The decision whether or not to publish a scholarly article, or what should be modified before publication, ultimately lies with the publisher (editor-in-chief or the editorial board) to which the manuscript has been submitted. Similarly, the decision whether or not to fund a proposed project rests with an official of the funding agency. These individuals usually refer to the opinion of one or more reviewers in making their decision. This is primarily for three reasons:
Workload. A small group of editors/assessors cannot devote sufficient time to each of the many articles submitted to many journals.
Miscellany of ideas. Were the editor/assessor to judge all submitted material themselves, approved material would solely reflect their opinion.
Limited expertise. An editor/assessor cannot be expected to be sufficiently expert in all areas covered by a single journal or funding agency to adequately judge all submitted material.
Reviewers are often anonymous and independent. However, some reviewers may choose to waive their anonymity, and in other limited circumstances, such as the examination of a formal complaint against the referee, or a court order, the reviewer's identity may have to be disclosed. Anonymity may be unilateral or reciprocal (single- or double-blinded reviewing).
Since reviewers are normally selected from experts in the fields discussed in the article, the process of peer review helps to keep some invalid or unsubstantiated claims out of the body of published research and knowledge. Scholars will read published articles outside their limited area of detailed expertise, and then rely, to some degree, on the peer-review process to have provided reliable and credible research that they can build upon for subsequent or related research. Significant scandal ensues when an author is found to have falsified the research included in an article, as other scholars, and the field of study itself, may have relied upon the invalid research.
For US universities, peer reviewing of books before publication is a requirement for full membership of the Association of American University Presses.
Procedure
In the case of proposed publications, the publisher (editor-in-chief or the editorial board, often with assistance of corresponding or associate editors) sends advance copies of an author's work or ideas to researchers or scholars who are experts in the field (known as "referees" or "reviewers"). Communication is nowadays normally by e-mail or through a web-based manuscript processing system such as ScholarOne, Scholastica, or Open Journal Systems. Depending on the field of study and on the specific journal, there are usually one to three referees for a given article. For example, Springer states that there are two or three reviewers per article.
The peer-review process involves three steps:
Step 1: Desk evaluation
An editor evaluates the manuscript to judge whether the paper will be passed on to journal referees. At this phase many articles receive a "desk reject", that is, the editor chooses not to pass along the article. The authors may or may not receive a letter of explanation.
Desk rejection is intended to be a streamlined process so that editors may move past nonviable manuscripts quickly and provide authors with the opportunity to pursue a more suitable journal. For example, the European Accounting Review editors subject each manuscript to three questions to decide whether a manuscript moves forward to referees: 1) Is the article a fit for the journal's aims and scope, 2) is the paper content (e.g. literature review, methods, conclusions) sufficient and does the paper make a worthwhile contribution to the larger body of literature, and 3) does it follow format and technical specifications? If "no" to any of these, the manuscript receives a desk rejection.
Desk rejection rates vary by journal. For example, in 2017 researchers at the World Bank compiled rejection rates of several global economics journals; the desk rejection rate ranged from 21% (Economic Lacea) to 66% (Journal of Development Economics). The American Psychological Association publishes rejection rates for several major publications in the field, and although they do not specify whether the rejection is pre- or post- desk evaluation, their figures in 2016 ranged from a low of 49% to a high of 90%.
Step 2: External review
If the paper is not desk rejected, the editors send the manuscript to the referees, who are chosen for their expertise and distance from the authors. At this point, referees may reject, accept without changes (rare) or instruct the authors to revise and resubmit.
Reasons vary for acceptance of an article by editors, but Elsevier published an article where three editors weigh in on factors that drive article acceptance. These factors include whether the manuscript: delivers "new insight into an important issue", will be useful to practitioners, advances or proposes a new theory, raises new questions, has appropriate methods and conclusion, presents a good argument based on the literature, and tells a good story. One editor notes that he likes papers that he "wished he'd done" himself.
These referees each return an evaluation of the work to the editor, noting weaknesses or problems along with suggestions for improvement. Typically, most of the referees' comments are eventually seen by the author, though a referee can also send 'for your eyes only' comments to the publisher; scientific journals observe this convention almost universally. The editor then evaluates the referees' comments, her or his own opinion of the manuscript before passing a decision back to the author(s), usually with the referees' comments.
Referees' evaluations usually include an explicit recommendation of what to do with the manuscript or proposal, often chosen from options provided by the journal or funding agency. For example, Nature recommends four courses of action:
to unconditionally accept the manuscript or the proposal,
to accept it in the event that its authors improve it in certain ways
to reject it, but encourage revision and invite re-submission
to reject it outright.
During this process, the role of the referees is advisory. The editor(s) is typically under no obligation to accept the opinions of the referees, though he or she will most often do so. Furthermore, the referees in scientific publication do not act as a group, do not communicate with each other, and typically are not aware of each other's identities or evaluations. Proponents argue that if the reviewers of a paper are unknown to each other, the editor(s) can more easily verify the objectivity of the reviews. There is usually no requirement that the referees achieve consensus, with the decision instead often made by the editor(s) based on her best judgement of the arguments.
In situations where multiple referees disagree substantially about the quality of a work, there are a number of strategies for reaching a decision. The paper may be rejected outright, or the editor may choose which reviewer's point the authors should address. When a publisher receives very positive and very negative reviews for the same manuscript, the editor will often solicit one or more additional reviews as a tie-breaker. As another strategy in the case of ties, the publisher may invite authors to reply to a referee's criticisms and permit a compelling rebuttal to break the tie. If a publisher does not feel confident to weigh the persuasiveness of a rebuttal, the publisher may solicit a response from the referee who made the original criticism. An editor may convey communications back and forth between authors and a referee, in effect allowing them to debate a point.
Even in these cases, however, publishers do not allow multiple referees to confer with each other, though each reviewer may often see earlier comments submitted by other reviewers. The goal of the process is explicitly not to reach consensus or to persuade anyone to change their opinions, but instead to provide material for an informed editorial decision. One early study regarding referee disagreement found that agreement was greater than chance, if not much greater than chance, on six of seven article attributes (e.g. literature review and final recommendation to publish), but this study was small and it was conducted on only one journal. At least one study has found that reviewer disagreement is not common, but this study is also small and on only one journal.
Traditionally, reviewers would often remain anonymous to the authors, but this standard varies both with time and with academic field. In some academic fields, most journals offer the reviewer the option of remaining anonymous or not, or a referee may opt to sign a review, thereby relinquishing anonymity. Published papers sometimes contain, in the acknowledgments section, thanks to anonymous or named referees who helped improve the paper. For example, Nature journals provide this option.
Sometimes authors may exclude certain reviewers: one study conducted on the Journal of Investigative Dermatology found that excluding reviewers doubled the chances of article acceptance. Some scholars are uncomfortable with this idea, arguing that it distorts the scientific process. Others argue that it protects against referees who are biased in some manner (e.g. professional rivalry, grudges). In some cases, authors can choose referees for their manuscripts. mSphere, an open-access journal in microbial science, has moved to this model. Editor-in-Chief Mike Imperiale says this process is designed to reduce the time it takes to review papers and permit the authors to choose the most appropriate reviewers. But a scandal in 2015 shows how this choosing reviewers can encourage fraudulent reviews. Fake reviews were submitted to the Journal of the Renin-Angiotensin-Aldosterone System in the names of author-recommended reviewers, causing the journal to eliminate this option.
Step 3: Revisions
If the manuscript has not been rejected during peer review, it returns to the authors for revisions. During this phase, the authors address the concerns raised by reviewers. William Stafford Noble offers ten rules for responding to reviewers. His rules include:
"Provide an overview, then quote the full set of reviews"
"Be polite and respectful of all reviewers"
"Accept the blame"
"Make the response self-contained"
"Respond to every point raised by the reviewer"
"Use typography to help the reviewer navigate your response"
"Whenever possible, begin your response to each comment with a direct answer to the point being raised"
"When possible, do what the reviewer asks"
"Be clear about what changed relative to the previous version"
"If necessary, write the response twice" (i.e. write a version for "venting" but then write a version the reviewers will see)
Recruiting referees
At a journal or book publisher, the task of picking reviewers typically falls to an editor. When a manuscript arrives, an editor solicits reviews from scholars or other experts who may or may not have already expressed a willingness to referee for that journal or book division. Granting agencies typically recruit a panel or committee of reviewers in advance of the arrival of applications.
Referees are supposed to inform the editor of any conflict of interests that might arise. Journals or individual editors may invite a manuscript's authors to name people whom they consider qualified to referee their work. For some journals this is a requirement of submission. Authors are sometimes also given the opportunity to name natural candidates who should be disqualified, in which case they may be asked to provide justification (typically expressed in terms of conflict of interest).
Editors solicit author input in selecting referees because academic writing typically is very specialized. Editors often oversee many specialties, and can not be experts in all of them. But after an editor selects referees from the pool of candidates, the editor typically is obliged not to disclose the referees' identities to the authors, and in scientific journals, to each other. Policies on such matters differ among academic disciplines.
One difficulty with respect to some manuscripts is that, there may be few scholars who truly qualify as experts, people who have themselves done work similar to that under review. This can frustrate the goals of reviewer anonymity and avoidance of conflicts of interest. Low-prestige or local journals and granting agencies that award little money are especially handicapped with regard to recruiting experts.
A potential hindrance in recruiting referees is that they are usually not paid, largely because doing so would itself create a conflict of interest. Also, reviewing takes time away from their main activities, such as his or her own research. To the would-be recruiter's advantage, most potential referees are authors themselves, or at least readers, who know that the publication system requires that experts donate their time. Serving as a referee can even be a condition of a grant, or professional association membership. In general, because of the explosion of the electronic information and the disproportionate increase in journal number versus the steady increase in the number of scientists has created a peer review crisis. The system currently in place is not responding to modern needs and will inevitably perish, unless radical reforms are made promptly. The academic system should revolutionize and establish strict peer review activity criteria essential for promotion and tenure, based on established universal metrics. That is, reward reviewers academically as it rewards researchers, which is currently not the case. All other incentives have failed.
Referees have the opportunity to prevent work that does not meet the standards of the field from being published, which is a position of some responsibility. Editors are at a special advantage in recruiting a scholar when they have overseen the publication of his or her work, or if the scholar is one who hopes to submit manuscripts to that editor's publishing entity in the future. Granting agencies, similarly, tend to seek referees among their present or former grantees.
Peerage of Science was an independent service and a community where reviewer recruitment happens via Open Engagement: authors submit their manuscript to the service where it is made accessible for any non-affiliated scientist, and 'validated users' choose themselves what they want to review. The motivation to participate as a peer reviewer comes from a reputation system where the quality of the reviewing work is judged and scored by other users, and contributes to user profiles. Peerage of Science does not charge any fees to scientists, and does not pay peer reviewers. Participating publishers however pay to use the service, gaining access to all ongoing processes and the opportunity to make publishing offers to the authors.
With independent peer review services the author usually retains the right to the work throughout the peer review process, and may choose the most appropriate journal to submit the work to. Peer review services may also provide advice or recommendations on most suitable journals for the work. Journals may still want to perform an independent peer review, without the potential conflict of interest that financial reimbursement may cause, or the risk that an author has contracted multiple peer review services but only presents the most favorable one.
An alternative or complementary system of performing peer review is for the author to pay for having it performed. Example of such service provider was Rubriq (2013-2017), that for each work assigned peer reviewers who were financially compensated for their efforts.
Different styles
Anonymous and attributed
For most scholarly publications, the identity of the reviewers is kept anonymised (also called "blind peer review"). The alternative, attributed peer review involves revealing the identities of the reviewers. Some reviewers choose to waive their right to anonymity, even when the journal's default format is blind peer review.
In anonymous peer review, reviewers are known to the journal editor or conference organiser but their names are not given to the article's author. In some cases, the author's identity can also be anonymised for the review process, with identifying information stripped from the document before review. The system is intended to reduce or eliminate bias.
Some experts proposed blind review procedures for reviewing controversial research topics.
In double-blind peer review, which has been fashioned by sociology journals in the 1950s and remains more common in the social sciences and humanities than in the natural sciences, the identity of the authors is concealed from the reviewers ("blinded"), and vice versa, lest the knowledge of authorship or concern about disapprobation from the author bias their review. Critics of the double-blind review process point out that, despite any editorial effort to ensure anonymity, the process often fails to do so, since certain approaches, methods, writing styles, notations, etc., point to a certain group of people in a research stream, and even to a particular person.
In many fields of "big science", the publicly available operation schedules of major equipments, such as telescopes or synchrotrons, would make the authors' names obvious to anyone who would care to look them up. Proponents of double-blind review argue that it performs no worse than single-blind, and that it generates a perception of fairness and equality in academic funding and publishing. Single-blind review is strongly dependent upon the goodwill of the participants, but no more so than double-blind review with easily identified authors.
As an alternative to single-blind and double-blind review, authors and reviewers are encouraged to declare their conflicts of interest when the names of authors and sometimes reviewers are known to the other. When conflicts are reported, the conflicting reviewer can be prohibited from reviewing and discussing the manuscript, or his or her review can instead be interpreted with the reported conflict in mind; the latter option is more often adopted when the conflict of interest is mild, such as a previous professional connection or a distant family relation. The incentive for reviewers to declare their conflicts of interest is a matter of professional ethics and individual integrity. Even when the reviews are not public, they are still a matter of record and the reviewer's credibility depends upon how they represent themselves among their peers. Some software engineering journals, such as the IEEE Transactions on Software Engineering, use non-blind reviews with reporting to editors of conflicts of interest by both authors and reviewers.
A more rigorous standard of accountability is known as an audit. Because reviewers are not paid, they cannot be expected to put as much time and effort into a review as an audit requires. Therefore, academic journals such as Science, organizations such as the American Geophysical Union, and agencies such as the National Institutes of Health and the National Science Foundation maintain and archive scientific data and methods in the event another researcher wishes to replicate or audit the research after publication.
The traditional anonymous peer review has been criticized for its lack of accountability, the possibility of abuse by reviewers or by those who manage the peer review process (that is, journal editors), its possible bias, and its inconsistency, alongside other flaws. Eugene Koonin, a senior investigator at the National Center for Biotechnology Information, asserts that the system has "well-known ills" and advocates "open peer review".
Open peer review
Pre- and post-publication peer review
The process of peer review is not restricted to the publication process managed by academic journals. In particular, some forms of peer review can occur before an article is submitted to a journal and/or after it is published by the journal.
Pre-publication peer review
Manuscripts are typically reviewed by colleagues before submission, and if the manuscript is uploaded to preprint servers, such as ArXiv, BioRxiv or SSRN, researchers can read and comment on the manuscript. The practice to upload to preprint servers, and the activity of discussion heavily depend on the field, and it allows an open pre-publication peer review. The advantage of this method is speed and transparency of the review process. Anyone can give feedback, typically in form of comments, and typically not anonymously. These comments are also public, and can be responded to, therefore author-reviewer communication is not restricted to the typical 2–4 rounds of exchanges in traditional publishing. The authors can incorporate comments from a wide range of people instead of feedback from the typically 3–4 reviewers. The disadvantage is that a far larger number of papers are presented to the community without any guarantee on quality.
Post-publication peer review
After a manuscript is published, the process of peer review continues as publications are read, known as post-publication peer review. Readers will often send letters to the editor of a journal, or correspond with the editor via an on-line journal club. In this way, all "peers" may offer review and critique of published literature. The introduction of the "epub ahead of print" practice in many journals has made possible the simultaneous publication of unsolicited letters to the editor together with the original paper in the print issue.
A variation on this theme is open peer commentary, in which commentaries from specialists are solicited on published articles and the authors are invited to respond. Journals using this process solicit and publish non-anonymous commentaries on the "target paper" together with the paper, and with original authors' reply as a matter of course. Open peer commentary was first implemented by the anthropologist Sol Tax, who founded the journal Current Anthropology in 1957. The journal Behavioral and Brain Sciences, published by Cambridge University Press, was founded by Stevan Harnad in 1978 and modeled on Current Anthropology's open peer commentary feature. Psycoloquy (1990–2002) was based on the same feature, but this time implemented online. Since 2016 open peer commentary is also provided by the journal Animal Sentience.
In addition to journals hosting their own articles' reviews, there are also external, independent websites dedicated to post-publication peer-review, such as PubPeer which allows anonymous commenting of published literature and pushes authors to answer these comments. It has been suggested that post-publication reviews from these sites should be editorially considered as well. The megajournals F1000Research and ScienceOpen publish openly both the identity of the reviewers and the reviewer's report alongside the article.
Some journals use post-publication peer review as formal review method, instead of pre-publication review. This was first introduced in 2001, by Atmospheric Chemistry and Physics (ACP). More recently F1000Research, Qeios, and ScienceOpen were launched as megajournals with post-publication review as formal review method. At ACP, F1000Research, and Qeios peer reviewers are formally invited, much like at pre-publication review journals. Articles that pass peer review at those three journals are included in external scholarly databases.
Social media and informal peer review
Recent research has called attention to the use of social media technologies and science blogs as a means of informal, post-publication peer review, as in the case of the #arseniclife (or GFAJ-1) controversy. In December 2010, an article published in Scienceexpress (the ahead-of-print version of Science) generated both excitement and skepticism, as its authorsled by NASA astrobiologist Felisa Wolfe-Simonclaimed to have discovered and cultured a certain bacteria that could replace phosphorus with arsenic in its physiological building blocks. At the time of the article's publication, NASA issued press statements suggesting that the finding would impact the search for extraterrestrial life, sparking excitement on Twitter under the hashtag #arseniclife, as well as criticism from fellow experts who voiced skepticism via their personal blogs. Ultimately, the controversy surrounding the article attracted media attention, and one of the most vocal scientific criticsRosemary Redfieldformally published in July 2012 regarding her and her colleagues' unsuccessful attempt to replicate the NASA scientists' original findings.
Researchers following the impact of the #arseniclife case on social media discussions and peer review processes concluded the following:
Our results indicate that interactive online communication technologies can enable members in the broader scientific community to perform the role of journal reviewers to legitimize scientific information after it has advanced through formal review channels. In addition, a variety of audiences can attend to scientific controversies through these technologies and observe an informal process of post-publication peer review. (p 946)
Result-blind peer review
Studies which report a positive or statistically significant result are far more likely to be published than ones which do not. A counter-measure to this positivity bias is to hide or make unavailable the results in the paper, making journal acceptance more like scientific grant agencies reviewing research proposals. Versions include:
Result-blind peer review or results blind peer review, first proposed 1966: Reviewers receive an edited version of the submitted paper which omits the results and conclusion section. In a two-stage version, a second round of reviews or editorial judgment is based on the full paper version, which was first proposed in 1977.
Conclusion-blind review, proposed by Robin Hanson in 2007 extends this further asking all authors to submit a positive and a negative version, and only after the journal has accepted the article authors reveal which is the real version.
Pre-accepted articles or outcome-unbiased journals or advance publication review or registered reports or prior to results submission or early acceptance extends study pre-registration to the point that journals accepted or reject papers based on the version of the paper written before the results or conclusions have been made (an enlarged study protocol), but instead describes the theoretical justification, experimental design, and statistical analysis. Only once the proposed hypothesis and methodology have been accepted by reviewers, the authors would collect the data or analyze previously collected data. A limited variant of a pre-accepted article was The Lancet's study protocol review from 1997 to 2015 reviewed and published randomized trial protocols with a guarantee that the eventual paper would at least be sent out to peer review rather than immediately rejected. For example, Nature Human Behaviour has adopted the registered report format, as it "shift[s] the emphasis from the results of research to the questions that guide the research and the methods used to answer them". The European Journal of Personality defines this format: "In a registered report, authors create a study proposal that includes theoretical and empirical background, research questions/hypotheses, and pilot data (if available). Upon submission, this proposal will then be reviewed prior to data collection, and if accepted, the paper resulting from this peer-reviewed procedure will be published, regardless of the study outcomes."
The following journals used result-blind peer review or pre-accepted articles:
The European Journal of Parapsychology, under Martin Johnson (who proposed a version of Registered Reports in 1974), began accepting papers based on submitted designs and then publishing them, from 1976 to 1993, and published 25 RRs total
The International Journal of Forecasting used opt-in result-blind peer review and pre-accepted articles from before 1986 through 1996/1997.
The journal Applied Psychological Measurement offered an opt-in "advance publication review" process from 1989 to 1996, ending use after only 5 papers were submitted.
The JAMA Internal Medicine found in a 2009 survey that 86% of its reviewers would be willing to work in a result-blind peer review process, and ran a pilot experiment with a two-stage result-blind peer review, showing the unblinded step benefited positive studies more than negatives. but the journal does not currently use result-blind peer review.
The Center for Open Science encourages using "Registered Reports" (pre-accepted articles) beginning in 2013. As of October 2017, ~80 journals offer Registered Reports in general, have had special issues of Registered Reports, or limited acceptance of Registered Reports (e.g. replications only) including AIMS Neuroscience, Cortex, Perspectives on Psychological Science, Social Psychology, & Comparative Political Studies
Comparative Political Studies published results of its pilot experiment of 19 submissions of which 3 were pre-accepted in 2016. the process worked well but submissions were weighted towards quantitative experimental designs, and reduced the amount of 'fishing' as submitters and reviewers focused on theoretical backing, substantive importance of results, with attention to the statistical power and implications of a null result, concluding that "we can clearly state that this form of review lead to papers that were of the highest quality. We would love to see a top journal adopt results-free review as a policy, at very least allowing results-free review as one among several standard submission options."
Extended peer review
Extended peer review is the process of including people and groups with experience beyond that of working academics in the processes of assuring the quality of research. If conducted systematically, this can lead to more reliable, or applicable, results than a peer review process conducted purely by academics.
Criticism
Scholarly peer review has been subject to several criticisms, and various proposals for reforming the system have been suggested over the years. Many studies have emphasized the problems inherent to the process of peer review. Moreover, Ragone et al., have shown that there is a low correlation between peer review outcomes and the future impact measured by citations.
Various biomedical editors in particular have expressed criticism of peer review. A Cochrane review found little empirical evidence that peer review ensures quality in biomedical research, while a second systematic review and meta-analysis found a need for evidence-based peer review in biomedicine given the paucity of assessment of the interventions designed to improve the process.
To an outsider, the anonymous, pre-publication peer review process is opaque. Certain journals are accused of not carrying out stringent peer review in order to more easily expand their customer base, particularly in journals where authors pay a fee before publication. Richard Smith, MD, former editor of the British Medical Journal, has claimed that peer review is "ineffective, largely a lottery, anti-innovatory, slow, expensive, wasteful of scientific time, inefficient, easily abused, prone to bias, unable to detect fraud and irrelevant; Several studies have shown that peer review is biased against the provincial and those from low- and middle-income countries; Many journals take months and even years to publish and the process wastes researchers' time. As for the cost, the Research Information Network estimated the global cost of peer review at £1.9 billion in 2008."
In addition, Australia's Innovative Research Universities group (a coalition of seven comprehensive universities committed to inclusive excellence in teaching, learning and research in Australia) has found that "peer review disadvantages researchers in their early careers, when they rely on competitive grants to cover their salaries, and when unsuccessful funding applications often mark the end of a research idea".
Peer review publication is a common requirement for academic tenure. This requirement has been criticised on cultural grounds. In 2011, University of British Columbia assistant law professor, Lorna McCue, argued that emphasis on peer review publication was culturally inappropriate as it did not recognize the importance of Indigenous oral traditions. In 2018, the British Columbia Human Rights Tribunal found that this complaint was not justified .
There is an ongoing discussion about a peer-review crisis. In 2022 Inside Higher Ed reported a serious shortage of scholars to review submitted articles and bigger structural problems amplified by the COVID-19 pandemic.
Low-end distinctions in articles understandable to all peers
Brezis and Birukou show that the peer review process is not working properly due to reviewer bias. Additionally, they underline that the ratings are not robust, e.g., changing reviewers can have a dramatic impact on the review results. Two main elements affect the bias in the peer process:
The first element is that referees display homophily in their taste and perception of innovative ideas. So reviewers who are developing conventional ideas will tend to give low grades to innovative projects, while reviewers who have developed innovative ideas tend, by homophily, to give higher grades to innovative projects.
The second element leading to a high variance in the peer review process is that reviewers are not investing the same amount of time to analyze the projects (or equivalently are not with the same abilities). Brezis and Biruku show that this heterogeneity among referees will lead to seriously affect the whole peer review process, and will lead to main arbitrariness in the results of the process.
Medical researcher John Ioannidis expands on this second point, arguing that since the exams and other tests that people pass on their way from "layman" to "expert" focus on answering the questions in time and in accordance with a list of answers, and not on making precise distinctions (the latter of which would be unrecognizable to experts of lower cognitive precision), there is as much individual variation in the ability to distinguish causation from correlation among "experts" as there is among "laymen". Ioannidis argues that as a result, scholarly peer review by many "experts" allows only articles that are understandable at a wide range of cognitive precision levels including very low ones to pass, biasing publications towards favoring articles that infer causation from correlation while mislabelling articles that make the distinction as "incompetent overestimation of one's ability" on the side of the authors because some of the reviewing "experts" are cognitively unable to distinguish the distinction from alleged rationalization of specific conclusions. It is argued by Ioannidis that this makes peer review a cause of selective publication of false research findings while stopping publication of rigorous criticism thereof, and that further post-publication review repeats the same bias by selectively retracting the few rigorous articles that may have made it through initial pre-publication peer review while letting the low-end ones that confuse correlation and causation remain in print.
Tendency to Discourage Innovative Projects
The peer process is also in use for projects acceptance. (For projects, the acceptance rates are small and are between 1% and 20%, with an average of 10%. In the European H2020 calls, the acceptance rate is 1.8%.) Peer review is more problematic when choosing the projects to be funded since innovative projects are not highly ranked in the existing peer-review process. The peer-review process leads to conformity, i.e., the selection of less controversial projects and papers. This may even influence the type of proposals scholars will propose, since scholars need to find financing for their research as discussed by Martin, 1997: "A common informal view is that it is easier to obtain funds for conventional projects. Those who are eager to get funding are not likely to propose radical or unorthodox projects. Since you don't know who the referees are going to be, it is best to assume that they are middle-of-the-road. Therefore, the middle-of-the-road application is safer".
Peer review and trust
Researchers have peer-reviewed manuscripts prior to publishing them in a variety of ways since the 18th century. The main goal of this practice is to improve the relevance and accuracy of scientific discussions. Even though experts often criticize peer review for a number of reasons, the process is still often considered the "gold standard" of science. Occasionally however, peer review approves studies that are later found to be wrong and rarely deceptive or fraudulent results are discovered prior to publication. Thus, there seems to be an element of discord between the ideology behind and the practice of peer review. By failing to effectively communicate that peer review is imperfect, the message conveyed to the wider public is that studies published in peer-reviewed journals are "true" and that peer review protects the literature from flawed science. A number of well-established criticisms exist of many elements of peer review. In the following we describe cases of the wider impact inappropriate peer review can have on public understanding of scientific literature.
Multiple examples across several areas of science find that scientists elevated the importance of peer review for research that was questionable or corrupted. For example, climate change deniers have published studies in the Energy and Environment journal, attempting to undermine the body of research that shows how human activity impacts the Earth's climate. Politicians in the United States who reject the established science of climate change have then cited this journal on several occasions in speeches and reports.
At times, peer review has been exposed as a process that was orchestrated for a preconceived outcome. The New York Times gained access to confidential peer review documents for studies sponsored by the National Football League (NFL) that were cited as scientific evidence that brain injuries do not cause long-term harm to its players. During the peer review process, the authors of the study stated that all NFL players were part of a study, a claim that the reporters found to be false by examining the database used for the research. Furthermore, The Times noted that the NFL sought to legitimize the studies" methods and conclusion by citing a "rigorous, confidential peer-review process" despite evidence that some peer reviewers seemed "desperate" to stop their publication. Recent research has also demonstrated that widespread industry funding for published medical research often goes undeclared and that such conflicts of interest are not appropriately addressed by peer review. Conflict of interest is less likely to be picked up in double-blinded reviews since the reviewer does not know the identity of the authors.
Another problem that peer review fails to catch is ghostwriting, a process by which companies draft articles for academics who then publish them in journals, sometimes with little or no changes. These studies can then be used for political, regulatory and marketing purposes. In 2010, the US Senate Finance Committee released a report that found this practice was widespread, that it corrupted the scientific literature and increased prescription rates. Ghostwritten articles have appeared in dozens of journals, involving professors at several universities.
Just as experts in a particular field have a better understanding of the value of papers published in their area, scientists are considered to have better grasp of the value of published papers than the general public and to see peer review as a human process, with human failings, and that "despite its limitations, we need it. It is all we have, and it is hard to imagine how we would get along without it". But these subtleties are lost on the general public, who are often misled into thinking that being published in a journal with peer review is the "gold standard" and can erroneously equate published research with the truth. Thus, more care must be taken over how peer review, and the results of peer-reviewed research, are communicated to non-specialist audiences; particularly during a time in which a range of technical changes and a deeper appreciation of the complexities of peer review are emerging. This will be needed as the scholarly publishing system has to confront wider issues such as retractions and replication or reproducibility "crises".
Views of peer review
Peer review is often considered integral to scientific discourse in one form or another. Its gatekeeping role is supposed to be necessary to maintain the quality of the scientific literature and avoid a risk of unreliable results, inability to separate signal from noise, and slow scientific progress.
Shortcomings of peer review have been met with calls for even stronger filtering and more gatekeeping. A common argument in favor of such initiatives is the belief that this filter is needed to maintain the integrity of the scientific literature.
Calls for more oversight have at least two implications that are counterintuitive of what is known to be true scholarship.
The belief that scholars are incapable of evaluating the quality of work on their own, that they are in need of a gatekeeper to inform them of what is good and what is not.
The belief that scholars need a "guardian" to make sure they are doing good work.
Others argue that authors most of all have a vested interest in the quality of a particular piece of work. Only the authors could have, as Feynman (1974) puts it, the "extra type of integrity that is beyond not lying, but bending over backwards to show how you're maybe wrong, that you ought to have when acting as a scientist." If anything, the current peer review process and academic system could penalize, or at least fail to incentivize, such integrity.
Instead, the credibility conferred by the "peer-reviewed" label could diminish what Feynman calls the culture of doubt necessary for science to operate a self-correcting, truth-seeking process. The effects of this can be seen in the ongoing replication crisis, hoaxes, and widespread outrage over the inefficacy of the current system. It's common to think that more oversight is the answer, as peer reviewers are not at all lacking in skepticism. But the issue is not the skepticism shared by the select few who determine whether an article passes through the filter. It is the validation, and accompanying lack of skepticism, that comes afterwards. Here again more oversight only adds to the impression that peer review ensures quality, thereby further diminishing the culture of doubt and counteracting the spirit of scientific inquiry.
Quality researcheven some of our most fundamental scientific discoveriesdates back centuries, long before peer review took its current form. Whatever peer review existed centuries ago, it took a different form than it does in modern times, without the influence of large, commercial publishing companies or a pervasive culture of publish or perish. Though in its initial conception it was often a laborious and time-consuming task, researchers took peer review on nonetheless, not out of obligation but out of duty to uphold the integrity of their own scholarship. They managed to do so, for the most part, without the aid of centralised journals, editors, or any formalised or institutionalised process whatsoever. Supporters of modern technology argue that it makes it possible to communicate instantaneously with scholars around the globe, make such scholarly exchanges easier, and restore peer review to a purer scholarly form, as a discourse in which researchers engage with one another to better clarify, understand, and communicate their insights.
Such modern technology includes posting results to preprint servers, preregistration of studies, open peer review, and other open science practices. In all these initiatives, the role of gatekeeping remains prominent, as if a necessary feature of all scholarly communication, but critics argue that a proper, real-world implementation could test and disprove this assumption; demonstrate researchers' desire for more that traditional journals can offer; show that researchers can be entrusted to perform their own quality control independent of journal-coupled review. Jon Tennant also argues that the outcry over the inefficiencies of traditional journals centers on their inability to provide rigorous enough scrutiny, and the outsourcing of critical thinking to a concealed and poorly-understood process. Thus, the assumption that journals and peer review are required to protect scientific integrity seems to undermine the very foundations of scholarly inquiry.
To test the hypothesis that filtering is indeed unnecessary to quality control, many of the traditional publication practices would need to be redesigned, editorial boards repurposed if not disbanded, and authors granted control over the peer review of their own work. Putting authors in charge of their own peer review is seen as serving a dual purpose. On one hand, it removes the conferral of quality within the traditional system, thus eliminating the prestige associated with the simple act of publishing. Perhaps paradoxically, the removal of this barrier might actually result in an increase of the quality of published work, as it eliminates the cachet of publishing for its own sake. On the other hand, readers know that there is no filter so they must interpret anything they read with a healthy dose of skepticism, thereby naturally restoring the culture of doubt to scientific practice.
In addition to concerns about the quality of work produced by well-meaning researchers, there are concerns that a truly open system would allow the literature to be populated with junk and propaganda by those with a vested interest in certain issues. A counterargument is that the conventional model of peer review diminishes the healthy skepticism that is a hallmark of scientific inquiry, and thus confers credibility upon subversive attempts to infiltrate the literature. Allowing such "junk" to be published could make individual articles less reliable but render the overall literature more robust by fostering a "culture of doubt".
Allegations of bias and suppression
The interposition of editors and reviewers between authors and readers may enable the intermediators to act as gatekeepers. Some sociologists of science argue that peer review makes the ability to publish susceptible to control by elites and to personal jealousy.
The peer review process may sometimes impede progress and may be biased against novelty. A linguistic analysis of review reports suggests that reviewers focus on rejecting the applications by searching for weak points, and not on finding the high-risk/high-gain groundbreaking ideas that may be in the proposal. Reviewers tend to be especially critical of conclusions that contradict their own views, and lenient towards those that match them. At the same time, established scientists are more likely than others to be sought out as referees, particularly by high-prestige journals/publishers. As a result, ideas that harmonize with the established experts' are more likely to see print and to appear in premier journals than are iconoclastic or revolutionary ones. This accords with Thomas Kuhn's well-known observations regarding scientific revolutions. A theoretical model has been established whose simulations imply that peer review and over-competitive research funding foster mainstream opinion to monopoly.
Criticisms of traditional anonymous peer review allege that it lacks accountability, can lead to abuse by reviewers, and may be biased and inconsistent.
There have also been suggestions of gender bias in peer review, with male authors being likely to receive more favorable treatment. However, a 2021 study found no evidence for such bias (and found that in some respects female authors were treated more favourably).
Exploitation of free work
Most academic publishers do not financially compensate reviewers for their participation in the peer-review process, which has been criticized by the academic community. Whereas some publishers have contended that it is economically not feasible to pay reviewers, some journals have started to pay reviewers through platforms such as Research Square when they are unable to receive free reviews. Other publishers such as Advances.in have made paying reviewers an inherent part of their business model.
Open access journals and peer review
Some critics of open access (OA) journals have argued that, compared to traditional subscription journals, open access journals might utilize substandard or less formal peer review practices, and, as a consequence, the quality of scientific work in such journals will suffer. In a study published in 2012, this hypothesis was tested by evaluating the relative "impact" (using citation counts) of articles published in open access and subscription journals, on the grounds that members of the scientific community would presumably be less likely to cite substandard work, and that citation counts could therefore act as one indicator of whether or not the journal format indeed impacted peer review and the quality of published scholarship. This study ultimately concluded that "OA journals indexed in Web of Science and/or Scopus are approaching the same scientific impact and quality as subscription journals, particularly in biomedicine and for journals funded by article processing charges," and the authors consequently argue that "there is no reason for authors not to choose to publish in OA journals just because of the 'OA' label.
Failures
Peer review fails when a peer-reviewed article contains fundamental errors that undermine at least one of its main conclusions and that could have been identified by more careful reviewers. Many journals have no procedure to deal with peer review failures beyond publishing letters to the editor. Peer review in scientific journals assumes that the article reviewed has been honestly prepared. The process occasionally detects fraud, but is not designed to do so. When peer review fails and a paper is published with fraudulent or otherwise irreproducible data, the paper may be retracted. A 1998 experiment on peer review with a fictitious manuscript found that peer reviewers failed to detect some manuscript errors and the majority of reviewers may not notice that the conclusions of the paper are unsupported by its results.
Fake peer review
There have been instances where peer review was claimed to be performed but in fact was not; this has been documented in some predatory open access journals (e.g., the Who's Afraid of Peer Review? affair) or in the case of sponsored Elsevier journals.
In November 2014, an article in Nature exposed that some academics were submitting fake contact details for recommended reviewers to journals, so that if the publisher contacted the recommended reviewer, they were the original author reviewing their own work under a fake name. The Committee on Publication Ethics issued a statement warning of the fraudulent practice. In March 2015, BioMed Central retracted 43 articles and Springer retracted 64 papers in 10 journals in August 2015. Tumor Biology journal is another example of peer review fraud.
In 2020, the Journal of Nanoparticle Research fell victim to an "organized rogue editor network", who impersonated respected academics, got a themed issue created, and got 19 substandard articles published (out of 80 submitted). The journal was praised for dealing with the scam openly and transparently.
Plagiarism
Reviewers generally lack access to raw data, but do see the full text of the manuscript, and are typically familiar with recent publications in the area. Thus, they are in a better position to detect plagiarism of prose than fraudulent data. A few cases of such textual plagiarism by historians, for instance, have been widely publicized.
On the scientific side, a poll of 3,247 scientists funded by the U.S. National Institutes of Health found 0.3% admitted faking data and 1.4% admitted plagiarism. Additionally, 4.7% of the same poll admitted to self-plagiarism or autoplagiarism, in which an author republishes the same material, data, or text, without citing their earlier work. Self-plagiarisms are less likely to be detected in double-blinded peer reviews.
Examples
"Perhaps the most widely recognized failure of peer review is its inability to ensure the identification of high-quality work. The list of important scientific papers that were rejected by some peer-reviewed journals goes back at least as far as the editor of Philosophical Transaction's 1796 rejection of Edward Jenner's report of the first vaccination against smallpox."
The Soon and Baliunas controversy involved the publication in 2003 of a review study written by aerospace engineer Willie Soon and astronomer Sallie Baliunas in the journal Climate Research, which was quickly taken up by the G.W. Bush administration as a basis for amending the first Environmental Protection Agency Report on the Environment. The paper was strongly criticized by numerous scientists for its methodology and for its misuse of data from previously published studies, prompting concerns about the peer review process of the paper. The controversy resulted in the resignation of several editors of the journal and the admission by its publisher Otto Kinne that the paper should not have been published as it was.
The trapezoidal rule, in which the method of Riemann sums for numerical integration was republished in a Diabetes research journal, Diabetes Care. The method is almost always taught in high school calculus, and was thus considered an example of an extremely well known idea being re-branded as a new discovery.
A conference organized by the Wessex Institute of Technology was the target of an exposé by three researchers who wrote nonsensical papers (including one that was composed of random phrases). They reported that the papers were "reviewed and provisionally accepted" and concluded that the conference was an attempt to "sell" publication possibilities to less experienced or naive researchers. This may however be better described as a lack of any actual peer review, rather than peer review having failed.
In the humanities, one of the most infamous cases of plagiarism undetected by peer review involved Martin Stone, formerly professor of medieval and Renaissance philosophy at the Hoger Instituut voor Wijsbegeerte of the KU Leuven. Martin Stone managed to publish at least forty articles and book chapters that were almost entirely stolen from the work of others. Most of these publications appeared in highly rated peer-reviewed journals and book series.
The controversial Younger Dryas impact hypothesis, which evolved directly from pseudoscience and now forms the basis for the pseudoarchaeology of Graham Hancock's Ancient Apocalypse, was first published in the peer-reviewed journal PNAS using a nonstandard review system, according to a comprehensive refutation by Holliday et al. (2023). According to this 2023 review, "Claiming evidence where none exists and providing misleading citations may be accidental, but when conducted repeatedly, it becomes negligent and undermines scientific advancement as well as the credibility of science itself. Also culpable is the failure of the peer review process to prevent such errors of fact from entering the literature. The Proceedings of the National Academy of Sciences 'contributed review' system for National Academy members...is at least partially responsible. The 'pal reviews' (as some refer to them) were significantly curtailed in 2010, in part due to the YDIH controversy."
Proposed Alternatives
Other attempts to reform the peer review process originate among others from the fields of metascience and journalology. Reformers seek to increase the reliability and efficiency of the peer review process and to provide it with a scientific foundation. Alternatives to common peer review practices have been put to the test, in particular open peer review, where the comments are visible to readers, generally with the identities of the peer reviewers disclosed as well, e.g., F1000, eLife, BMJ, and BioMed Central. In the case of eLife, peer review is used not for deciding whether to publish an article, but for assessing its importance and reliability. Likewise, the recognition and recruitment of peer reviewers continues to be a significant issue in the field of scholarly publishing.
In popular culture
In 2017, the Higher School of Economics in Moscow unveiled a "Monument to an Anonymous Peer Reviewer". It takes the form of a large concrete cube, or dice, with "Accept", "Minor Changes", "Major Changes", "Revise and Resubmit" and "Reject" on its five visible sides. Sociologist Igor Chirikov, who devised the monument, said that while researchers have a love-hate relationship with peer review, peer reviewers nonetheless do valuable but mostly invisible work, and the monument is a tribute to them.
| Physical sciences | Science basics | Basics and measurement |
3895284 | https://en.wikipedia.org/wiki/Tidal%20range | Tidal range | Tidal range is the difference in height between high tide and low tide. Tides are the rise and fall of sea levels caused by gravitational forces exerted by the Moon and Sun, by Earth's rotation and by centrifugal force caused by Earth's progression around the Earth-Moon barycenter. Tidal range depends on time and location.
Larger tidal range occur during spring tides (spring range), when the gravitational forces of both the Moon and Sun are aligned (at syzygy), reinforcing each other in the same direction (new moon) or in opposite directions (full moon). The largest annual tidal range can be expected around the time of the equinox if it coincides with a spring tide. Spring tides occur at the second and fourth (last) quarters of the lunar phases.
By contrast, during neap tides, when the Moon and Sun's gravitational force vectors act in quadrature (making a right angle to the Earth's orbit), the difference between high and low tides (neap range) is smallest. Neap tides occur at the first and third quarters of the lunar phases.
Tidal data for coastal areas is published by national hydrographic offices. The data is based on astronomical phenomena and is predictable. Sustained storm-force winds blowing from one direction combined with low barometric pressure can increase the tidal range, particularly in narrow bays. Such weather-related effects on the tide can cause ranges in excess of predicted values and can cause localized flooding. These weather-related effects are not calculable in advance.
Mean tidal range is calculated as the difference between mean high water (i.e., the average high tide level) and mean low water (the average low tide level).
Geography
The typical tidal range in the open ocean is about – mapped in blue and green at right. Mean ranges near coasts vary from near zero to , with the range depending on the volume of water adjacent to the coast, and the geography of the basin the water sits in. Larger bodies of water have higher ranges, and the geography can act as a funnel amplifying or dispersing the tide.
The world's largest mean tidal range of occurs in the Bay of Fundy, Canada (more specificially, at Burntcoat Head, Nova Scotia). The next highest, of , is at Ungava Bay, also in Canada, and the next, of , in the Bristol Channel, between England and Wales. The highest predicted extreme (not mean) range is , in the Bay of Fundy. The maximum range in the Bristol Channel is . The fifty coastal locations with the largest ranges worldwide are listed by the National Oceanic and Atmospheric Administration of the United States.
Some of the smallest tidal ranges occur in the Mediterranean, Baltic, and Caribbean Seas. A point within a tidal system where the tidal range is almost zero is called an amphidromic point.
Classification
The tidal range has been classified as:
Micro-tidal – when the tidal range is lower than 2 metres (6'6¾").
Meso-tidal – when the tidal range is between 2 metres and 4 metres (6'6¾" and 13'1½").
Macro-tidal – when the tidal range is higher than 4 metres (13'1½").
| Physical sciences | Oceanography | Earth science |
3898548 | https://en.wikipedia.org/wiki/Borophagus | Borophagus | Borophagus ("gluttonous eater") is an extinct genus of the subfamily Borophaginae, a group of canids endemic to North America from the Middle Miocene epoch through the Early Pleistocene epoch 12—1.8 Mya.
Evolution
Borophagus, like other borophagines, are loosely known as "bone-crushing" or "hyena-like" dogs. Though not the most massive borophagine by size or weight, it had a more highly evolved capacity to crunch bone than earlier, larger genera such as Epicyon, which seems to be an evolutionary trend of the group (Turner, 2004). During the Pliocene epoch, Borophagus began being displaced by other Canid species such as Canis edwardii and later by Aenocyon dirus. Early species of Borophagus were placed in the genus Osteoborus until recently, but the genera are now considered synonyms.
Description
Typical features of this genus are a bulging forehead and powerful jaws; Borophagus has been considered to be probably a scavenger by paleontologists in the past. Its crushing premolar teeth and strong jaw muscles would have been used to crack open bone, much like the hyena of the Old World. However, Borophagus fossils are so abundant and geographically widespread that some paleontologists now argue that Borophagus must have been both the dominant carnivore of its time, and thus an active predator because carrion feeding alone could not have sustained such a large population. They note that not all carnivores with bone-cracking ability are scavengers, such as the modern spotted hyena; instead, they interpret the bone-cracking ability as an adaptation to social hunting where complete utilization of a carcass was favored. Coprolites from Borophagus further vindicate its bone-crushing abilities, while simultaneously indicating it occupied a niche no longer seen in the present-day ecosystems of North America. The discovery of these coprolites also indicates that Borophagus may have been a social pack-hunter.
The adult animal is estimated to have been about 80 cm in length, similar to a coyote, although it was much more powerfully built.
Species
Borophagus diversidens existed for (synonymous with Felis hillianus, Hyaenognathus matthewi, Hyaenognathus pachyodon, Hyaenognathus solus, Porthocyon dubius)
Borophagus dudleyi existed for
Borophagus hilli existed for (synonymous with Osteoborus crassapineatus, Osteoborus progressus)
Borophagus littoralis existed for (syn. Osteoborus diabloensis)
Borophagus orc existed for
Borophagus parvus existed for
Borophagus pugnator existed for (synonymous with Osteoborus galushai)
Borophagus secundus existed for (synonymous with Hyaenognathus cyonoides, Hyaenognathus direptor, Osteoborus secundus)
Existence based on Figure 141 of Wang et al. (1999).
Paleoecology
In North America, in places such as Coffee Ranch in Texas, Borophagus was contemporary with the bear Agriotherium as well as the feliform Barbourofelis, the saber-toothed machairodont cat Amphimachairodus coloradensis and fellow canid Epicyon. All of these animals were potential competitors that would have occasionally conflicted with Borophagus for food and territory, though it may also have readily scavenged their kills. Prey for Borophagus included herbivores like the camel Aepycamelus, the pronghorn antelope Cosoryx, horses like Neohipparion and Nannippus, the ancient peccary Prosthennops and even rhinoceroses like the hippo-like Teleoceras, all of which could provide a suitable meal through hunting or scavenging.
| Biology and health sciences | Other carnivora | Animals |
27815938 | https://en.wikipedia.org/wiki/Parent%20hydride | Parent hydride | In chemistry, a parent hydride in IUPAC nomenclature refers to a main group compound with the formula , where A is a main group element. The names of parent hydrides end with -ane, analogous with the nomenclature for alkanes. Derivatives of parent hydrides are named by appending prefixes or suffixes to the name of the parent hydride to indicate the substituents that replace the hydrogen atoms.
Parent hydrides are used in both the organic nomenclature, and inorganic nomenclature systems.
*extensive body of chemistry
Reactions and structure
Parent hydrides are useful reference compounds, but many are nonexistent or unstable. Group III parent hydrides exist only under extraordinary conditions. Borane dimerizes irreversibly. Gallane and heavier congeners polymerize, sometimes with loss of hydrogen. Plumbane and bismuthane, stibane, indane, thallane are unstable.
| Physical sciences | Nomenclature | Chemistry |
24521450 | https://en.wikipedia.org/wiki/Diaphoretickes | Diaphoretickes | Diaphoretickes ( ) is a major group of eukaryotic organisms, with over 400,000 species. The majority of the earth's biomass that carries out photosynthesis belongs to Diaphoretickes.
Phylogeny
In 2012 the term Diaphoretickes was coined, and received the following phylogenetic definition:
The placement of Hemimastigophora is uncertain. Some studies find it nested within the clades of Diaphoretickes, as the best supported hypothesis or a less-supported one. Others find it to be the sister group of all other Diaphoretickes, which under the branch-based definition falls within Diaphoretickes.
Diaphoretickes includes the following groups:
In 2024, Hemimastigophora and Provora were placed as sister to Diaphoretickes, together with Meteora.
Taxonomic history
Diaphoretickes was identified by Burki et al. (2008) as the
"plants+HC+SAR megagroup". because it included plants (Archaeplastida), haptophytes, cryptomonads, and stramenopiles, alveolates, and rhizarians. Diaphoretickes has been called the SAR/HA Supergroup or "Corticata with Rhizaria". According to this description, it includes most of the species engaging in photosynthesis, except for the Euglenozoa and Cyanobacteria. It includes all Bikonts that are not excavates. The name "Corticata" comes from Cavalier-Smith's hypothesis about the common origin of the cortical alveoli of glaucophytes and alveolates. The megagroup was previously described as the sum of Archaeplastida, Rhizaria, and chromalveolates. However, this description is obsolete, largely due to the discovery that Chromalveolata, a refinement of Chromista, is not monophyletic.
| Biology and health sciences | Bikonts | Plants |
2908792 | https://en.wikipedia.org/wiki/Acasta%20Gneiss | Acasta Gneiss | The Acasta Gneiss Complex, also called the Acasta Gneiss, is a body of felsic to ultramafic Archean basement rocks, gneisses, that form the northwestern edge of the Slave Craton in the Northwest Territories, Canada, about north of Yellowknife, Canada. This geologic complex consists largely of tonalitic and granodioritic gneisses and lesser amounts of mafic and ultramafic gneisses. It underlies and is largely concealed by thin, patchy cover of Quaternary glacial sediments over an area of about . The Acasta Gneiss Complex contains fragments of the oldest known crust and record of more than a billion years (>4.0–2.9 Ga) of magmatism and metamorphism. The Acasta Gneiss Complex is exposed in a set of anticlinoriums within the foreland fold and thrust belt of the Paleoproterozoic Wopmay Orogen.
Nomenclature
The gneisses of the Acasta Gneiss Complex were first regionally mapped during a mapping program carried out by the Geological Survey of Canada. These gneisses were informally named the Acasta gneiss complex for the Acasta River east of Great Bear Lake in 1985 by J. E. King. Later, it was referred to as the Acasta gneisses by S.A. Bowring and others in 1989 In subsequent publications, Acasta Gneiss Complex is used instead of Acasta gneisses.
As the result of a detailed mapping and U–Pb dating, the oldest zircons of the Acasta Gneiss Complex were found to occur within a relatively homogeneous and mappable tonalitic gneiss by J. R. Reimink and others in 2014. They named it the Idiwhaa tonalitic gneiss. The Idiwhaa tonalitic gneiss has yielded abundant well-preserved igneous zircons with a U–Pb crystallization age of 4.02 Ga. It is a relatively homogeneous mafic tonalite that consists primarily of plagioclase, quartz, hornblende, biotite and minor garnet, with small, cross-cutting felsic veins.
Lithology
The Acasta Gneiss Complex is a heterogeneous assemblage of foliated to gneissic tonalites, trondhjemites, granodiorites, and granites which contains minor quartz-diorites, diorites, gabbros, and ultramafic rocks. Of the of the Canadian Shield underlain by Acasta Gneiss, only a area has been mapped and studied great detail of its remoteness and patchy outcrops due to a thick cover of Pleistocene glacial deposits. The study of the age distribution of detrital zircons from the Pleistocene deposits inferred the unmapped Acasta Gneiss Complex covered by Pleistocene deposits consists of a significant volume of 3.37 Ga granitoids. This study also concluded that near its inferred eastern limit, the unmapped Acasta Gneiss Complex contain rocks at least as old as 3.95 Ga.
Contacts
The base of Acasta Gneiss Complex is unexposed and unknown. It is considered to be unconformibly overlain by Central Slave Cover Group. The base of younger Late Mesoarchean greenstone belts that comprise the Central Slave Cover Group consists largely of fuchsite-bearing quartz arenite,. These quartz arenites are associated with chert-magnetite iron formation, rhyolite, and locally some ultrabasic rocks. These strata are intrerpreted to be remnants of the earliest sedimentary and volcanic strata that accumulated on Mesoarchean crystalline basement of the Slav craton.<ref name="IsachsenOthers1994a">Isachsen, C.E. and Bowring, S.A., 1994. Evolution of the Slave craton. Geology, 22(10), pp. 917-920.</ref>
Origin
From decades of petrologic, geochronologic, and geochemical studies, it has become apparent that the Acasta Gneiss Complex, as it is now exposed in the Canadian Shield, is the result of the complex interplay between repeated periods of felsic magma generation from the partial melting of older, mafic crust, its crystallization as felsic plutonic rocks, and subsequent thermal metamorphism.Bowring S.A. and Williams I.S., 1999. Priscoan (4.00-4.03 Ga) orthogneisses from northwestern Canada. Contributions to Mineralogy and Petrology, 134, p. 3-16. The formation of the Acasta Gneiss Complex began at the end of the Hadean and continued episodically over much of the Archean. The oldest felsic rocks that are currently exposed intruded pre-existing, even older, mafic crustal rock and crystallized well beneath the Earth's surface. Later, these rocks were thermally metamorphosed, intruded by additional felsic magma, and partially melted during Eoarchean thermal events occurred about 3.85 to 3.72 Ga and 3.66 to 3.59 Ga. The metamorphism and intrusion by felsic magmas continued throughout the Archean. This included major periods of intrusion by granitic magmas about 3.30 Ga, 2.88 Ga, and 2.70 Ga. Finally, these rocks were intruded by multiple Paleoproterozoic dike swarms and syenites about 1.80 Ga. The setting in which the tectonic activity that created the Acasta Gneiss Complex is proposed to have been a block of dominatly mafic protocrust, at least some of which crystallized in the Hadean. This block of protocrust would have been very similar to modern Iceland in its geology.
Earth’s oldest rock
With zircon U-Pb crystallization ages as old as about 4.03 Ga, beds within the Acasta Gneiss Complex are, , regarded to be the oldest known felsic rocks (crust) on Earth. Based on the analyses of Sm and Nd isotopes, mafic rocks of the Nuvvuagittuq Greenstone Belt, found on the shore of Hudson Bay, Inukjuak, Quebec, are argued to be the oldest rocks (crust) on Earth with estimated crystallization age of 4.313 Ga. Unfortunately, because of the basaltic composition and low-Zr content of the rocks comprising Nuvvuagittuq Greenstone Belt, they essentially lack the igneous zircons needed to reliably date them. As a result, their age and status as oldest known rock on Earth remains controversial. U-Pb crystallization ages of igneous zircons from felsic, trondhjemitic bands of intrusive origin within the Nuvvuagittuq Greenstone Belt only demonstrate that these rocks are older than 3.768 Ga.
Exhibit
In 2003 a team from the Smithsonian Institution collected a four-ton boulder of Acasta Gneiss Complex for display outside the National Museum of the American Indian in Washington, D.C. Another sample is on display in the Museu de Geociências of the University of Brasília, Brazil.
In October 2016, a sample of Acasta Gneiss Complex was placed on public display at the Clark Planetarium in Salt Lake City, Utah, US alongside samples of stromatolite and banded iron formation.
In 2006, Peter Skinner and Bert Cervo contributed a small sample of the rock to the Six String Nation project. Part of that material was inlaid into the first fret of Voyageur'', the guitar at the heart of the project.
On the campus of the University of Waterloo in Waterloo, Canada, a boulder from the Acasta Gneiss Complex has been placed in the Peter Russell Rock Garden.
| Physical sciences | Geologic features | Earth science |
2909308 | https://en.wikipedia.org/wiki/Pressure%20gradient | Pressure gradient | In hydrodynamics and hydrostatics, the pressure gradient (typically of air but more generally of any fluid) is a physical quantity that describes in which direction and at what rate the pressure increases the most rapidly around a particular location. The pressure gradient is a dimensional quantity expressed in units of pascals per metre (Pa/m). Mathematically, it is the gradient of pressure as a function of position. The gradient of pressure in hydrostatics is equal to the body force density (generalised Stevin's Law).
In petroleum geology and the petrochemical sciences pertaining to oil wells, and more specifically within hydrostatics, pressure gradients refer to the gradient of vertical pressure in a column of fluid within a wellbore and are generally expressed in pounds per square inch per foot (psi/ft). This column of fluid is subject to the compound pressure gradient of the overlying fluids. The path and geometry of the column is totally irrelevant; only the vertical depth of the column has any relevance to the vertical pressure of any point within its column and the pressure gradient for any given true vertical depth.
Physical interpretation
The concept of a pressure gradient is a local characterisation of the air (more generally of the fluid under investigation). The pressure gradient is defined only at these spatial scales at which pressure (more generally fluid dynamics) itself is defined.
Within planetary atmospheres (including the Earth's), the pressure gradient is a vector pointing roughly downwards, because the pressure changes most rapidly vertically, increasing downwards (see vertical pressure variation). The value of the strength (or norm) of the pressure gradient in the troposphere is typically of the order of 9 Pa/m (or 90 hPa/km).
The pressure gradient often has a small but critical horizontal component, which is largely responsible for wind circulation in the atmosphere. The horizontal pressure gradient is a two-dimensional vector resulting from the projection of the pressure gradient onto a local horizontal plane. Near the Earth's surface, this horizontal pressure gradient force is directed from higher toward lower pressure. Its particular orientation at any one time and place depends strongly on the weather situation. At mid-latitudes, the typical horizontal pressure gradient may take on values of the order of 10−2 Pa/m (or 10 Pa/km), although rather higher values occur within meteorological fronts.
Weather and climate relevance
Interpreting differences in air pressure between different locations is a fundamental component of many meteorological and climatological disciplines, including weather forecasting. As indicated above, the pressure gradient constitutes one of the main forces acting on the air to make it move as wind. Note that the pressure gradient force points from high towards low pressure zones. It is thus oriented in the opposite direction from the pressure gradient itself.
In acoustics
In acoustics, the pressure gradient is proportional to the sound particle acceleration according to Euler's equation. Sound waves and shock waves can induce very large pressure gradients, but these are oscillatory, and often transitory disturbances.
| Physical sciences | Thermodynamics | Physics |
37205436 | https://en.wikipedia.org/wiki/Saltwater%20fish | Saltwater fish | Saltwater fish, also called marine fish or sea fish, are fish that live in seawater. Saltwater fish can swim and live alone or in a large group called a school.
Saltwater fish are very commonly kept in aquariums for entertainment. Many saltwater fish are also caught to be eaten, or grown in aquaculture. However, many fish species have been overfished and are otherwise threatened by marine pollution or ecological changes caused by climate change.
Diet
Fishes that live in the ocean can be carnivores, herbivores, or omnivores. Herbivores in the ocean eat things such as algae and flowering seagrasses. Many herbivores' diets consist of primarily algae. Most saltwater fish will eat both macroalgae and microalgae. Many fish eat red, green, brown, and blue algae, but some fish prefer other types. Most saltwater fish that are carnivores will never eat algae under any circumstances. Carnivores' diets consist of shrimp, plankton, or tiny crustaceans.
Captivity
Saltwater aquariums are a multi-million dollar industry in the United States. About 10 million marine fish are imported into the United States each year for aquarium use. The United States imports more saltwater fish than any other country in the world. There are approximately 2,000 different species of saltwater fish that are imported and used in captivity. In many circumstances, fish used for marine trade are collected using harmful tactics such as cyanide. One way that people are trying to protect the coral reefs is by breeding marine fish in captivity. Captive-bred fish are known to be healthier and likely to live longer. Captive-bred fish are less susceptible to disease because they have not been exposed to the wild and they have not been damaged during the shipment process. Fish that are bred in captivity are already accustomed to aquarium habitats and food.
Habitats
There are many different components that make up a marine life habitat. Some of them are the temperature of the water, the quality, and quantity of water (flow and depth). Other components that can also contribute to the habitat of saltwater fish are pH level, salt level, and alkalinity level. Levels of nitrates and phosphates are also relevant, particularly when considering conditions for fish in captivity. There are other physical features that contribute to a habitat which are physical materials like rocks, reefs, and sand or the vegetation like the amount of algae, water plants, and saltmarsh. Specific fish live in specific habitats based on what they eat or what cycle of life they are currently at, another thing is the amount of salt that is in the water at that specific location. Some ocean habitats are not technically in the ocean and these are called estuaries, areas when oceans and rivers meet creating a mixture of salt water and freshwater making a different habitat for different types of fish and creatures to live in. The ocean is home to organisms as large as whales and as small as microscopic marine organisms such as phytoplankton. However, the vast majority of ocean life that humans are exposed to is simple saltwater fish. Saltwater fish can live in the deepest depths of the ocean where no sunlight can penetrate, but they can also live on the surface of the water.
Threats
Marine fish face many anthropogenic threats. Common human-induced threats include overfishing, pollution, habitat loss and destruction, climate change and invasive species. The aforementioned threats all come with a multitude of negative direct and indirect affects to marine ecosystems. With the human population growing at an exponential rate, these threats are likely to continue to be prevalent in marine ecosystems.
Overfishing
Overfishing is defined as the mass removal of fish from a body of water that results in halting the ability for breeding populations to replenish what has been removed. Fish is one of the most popular foods in the world and consumption has continued to rise with the growing human population and will continue to do so. The value of the global seafood market has seen a 15% increase from 2016 to 2020 and is projected to increase even more by 2023. Although it provides many people with a source of food, the global seafood market is a major threat to the biodiversity of fishes. Bycatch is a direct effect of overfishing and is defined as the unwanted capture of different marine organisms during industrial fishing. This results in many different species of fish dying after they are captured and discarded. Data on bycatch is often unclear and not well recorded but it is estimated that the U.S. alone discards 17–22% of their catch annually. The Mesopredator release hypothesis is one of the indirect effects of overfishing that is also often referred to as "fishing down the food web". This phenomenon means as fisheries deplete large apex predatory species, mid-sized predatory species increase in abundance and assume the role as top predators on the food web. This impacts the food web in marine environments and disrupts the balance of the ecosystem and is likely to cause trophic cascades.
Species affected by overfishing
Profitable fish stocks like the bluefin tuna are decreasing in numbers because of high demand. According to the IUCN Red List, the Pacific bluefin tuna, Atlantic bluefin tuna, and southern bluefin tuna are classified as vulnerable, endangered and critically endangered all due to over-exploitation.
According to the IUCN Red List, the oceanic whitetip shark is considered critically endangered because of its value in the seafood market. Their rapidly declining population is due to people overfishing them for their fins. They are a popular species of shark used in shark fin soup because of the size of their fins. All sharks are used for shark fin soup, however, certain species of sharks are preferred over others because of the large size of their fins.
The great white shark is listed as vulnerable on the IUCN Red List because its fins are commonly used in shark fin soup and has led to people over harvesting them for their fins. This shark belongs to the class Chondrichthyes that includes all sharks, skates and rays. The great white shark is one of the many examples of shark species threatened by human consumption because of the shark fin soup, large population declines of this class has been noted since the early 2000s due to the high demand for their fins, gill rakers and liver oil.
The Atlantic cod was historically abundant in the waters off the coast of New England. Due to its low fat content and dense white flesh, this fish is a popular choice among humans. Now considered vulnerable, its populations have both decreased in abundance and their distribution has shifted from northern to southern areas due to overfishing.
Cage aquaculture
Aquaculture is defined as the farming of aquatic organisms in controlled environments for the purpose of providing food and resources for humans. Aquaculture can take place in both marine and freshwater environments, however, because this is the saltwater fish page this entry will only cover the effects of aquaculture on marine fishes. The rising global demand for fish has contributed to the increase in aquaculture. Due to the decline of many wild fish stocks, aquaculture is the fastest growing food production system that contributes about 50% of the worlds fish supply. There is a lot of debate on whether or not aquaculture is an environmentally sustainable practice, yet the socioeconomic benefits that humans receive is tough to argue against. That being said there are significant negative effects that aquaculture, especially cage aquaculture, has on the surrounding environment.
Cage aquaculture involves rearing aquatic organisms in natural water sources while enclosed in a mesh/net cage that allows water from the surrounding environment to freely flow in and out. Cage aquaculture in marine environments has been particularly controversial because of the effects it has on the surrounding ecosystem thus, affecting wild marine fish populations. The main impacts of cage aquaculture are reduced water quality from fish sewage, high potential of genetic pollution of wild stocks due to escapees from aquaculture cages and the possibility of introducing an invasive species if the fish being reared are non-native . Fish sewage is the combination of fish feed, fecal material and antibiotics that is accumulated on the seafloor and in the water column from fish that are being farmed. It is not only harmful to wild fish stocks but it also poses a threat to marine plant life which is often a food source for wild fish stocks. Fish sewage is harmful because it pollutes the surrounding ecosystem and can cause problems like eutrophication, transmission of parasites and diseases to wild populations and developmental abnormalities on surrounding wild fish. Genetic pollution of wild fish populations is a common risk that cage aquaculture faces. For example, there are many scientific papers that have examined the effects of Atlantic Salmon escaping from their enclosures and interacting with wild populations. Farmed salmon have lower fitness (low survival rates and reproductive success) than a wild salmon would due to differences in artificial and natural selection. Artificial selection that chooses phenotypic traits that are desired for human consumption will alter the genetics of wild stocks if farmed fish interact and breed with wild populations. This would result in the reduction of fitness related traits that wild stocks possess which is a serious threat to these populations.
Categorization of saltwater fish by habitats
Coastal fish (also offshore fish or neritic fish) inhabit the sea between the shoreline and the edge of the continental shelf
Deep sea fish live below the photic zone of the ocean, i.e. where not enough light penetrates for photosynthesis to occur
Pelagic fish live near the surface of the sea or a lake
Demersal fish live on or near the bottom of the sea or a lake
Coral reef fish are associated with a coral reef.
| Biology and health sciences | Fishes by habitat | Animals |
1442361 | https://en.wikipedia.org/wiki/Four-bar%20linkage | Four-bar linkage | In the study of mechanisms, a four-bar linkage, also called a four-bar, is the simplest closed-chain movable linkage. It consists of four bodies, called bars or links, connected in a loop by four joints. Generally, the joints are configured so the links move in parallel planes, and the assembly is called a planar four-bar linkage. Spherical and spatial four-bar linkages also exist and are used in practice.
Planar four-bar linkage
Planar four-bar linkages are constructed from four links connected in a loop by four one-degree-of-freedom joints. A joint may be either a revolute joint – also known as a pin joint or hinged joint – denoted by R, or a prismatic joint – also known as a sliding pair – denoted by P.
A link that is fixed in place relative to the viewer is called a ground link.
A link connecting to the ground by a revolute joint that can perform a complete revolution is called a crank link.
A link connecting to the ground by a revolute joint that cannot perform a complete revolution is called a rocker link.
A link connecting to a ground line by a prismatic joint is called a slider. Sliders are sometimes considered to be cranks that have a hinged pivot at an infinitely long distance away perpendicular to the travel of the slider.
A link connecting to two other links is called a floating link or coupler.
A coupler connecting a crank and a slider in a single slider crank mechanism is often called a connecting rod, however, it has also been used to refer to any type of coupler.
There are three basic types of planar four-bar linkage, depending on the use of revolute or prismatic joints:
Four revolute joints:It is denoted as RRRR, constructed from four links connected by four revolute joints. The planar quadrilateral linkage refers to all arrangements in this type.Examples of 4R linkages include:
Double-crank linkage
Crank-rocker (Treadle) linkage (used in pumpjacks)
Double-rocker linkage (used in Ackermann steering)
Parallelogram (Parallel Motion) and Antiparallelogram (Contraparallelogram, Inverse Parallelogram, Butterfly, Bow-tie) linkages
Deltoid (Galloway) and Trapezium (Arglin) linkages
Three revolute joints:It is denoted as RRRP, PRRR, RPRR, or RRPR, constructed from four links connected by three revolute joints and one prismatic joint. The slider-crank linkage (RRRP) is one type of arrangement such that one link is a crank, which is then connected to a slider by a connecting rod. The inverted slider-crank is another type of arrangement such that there are two cranks with a slider acting as the coupler.Examples of 3R1P linkages include:
Single-slider crank mechanism (used in internal combustion engines)
Whitworth Quick Return mechanism (used in early types of shapers)
Crank and slotted lever Quick Return mechanism (used in shapers)
Fixed piston mechanism (used in hand pumps)
Two revolute joints and two prismatic joints:It is denoted as PRRP, and is constructed by connected two sliders with a coupler link. The doubler slider refers to all arrangements in this type.Examples of 2R2P linkages include:
Trammel of Archimedes (Elliptical trammel)
Scotch yoke (slotted link) mechanism (used in valve actuators)
Oldham's coupling
Planar four-bar linkages can be designed to guide a wide variety of movements, and are often the base mechanisms found in many machines. Because of this, the kinematics and dynamics of planar four-bar linkages are also important topics in mechanical engineering.
Planar quadrilateral linkage
Planar quadrilateral linkage, RRRR or 4R linkages have four rotating joints. One link of the chain is usually fixed, and is called the ground link, fixed link, or the frame. The two links connected to the frame are called the grounded links and are generally the input and output links of the system, sometimes called the input link and output link. The last link is the floating link, which is also called a coupler or connecting rod because it connects an input to the output.
Assuming the frame is horizontal there are four possibilities for the input and output links:
A crank: can rotate a full 360 degrees
A rocker: can rotate through a limited range of angles which does not include 0° or 180°
A 0-rocker: can rotate through a limited range of angles which includes 0° but not 180°
A π-rocker: can rotate through a limited range of angles which includes 180° but not 0°
Some authors do not distinguish between the types of rocker.
Grashof condition
The Grashof condition for a four-bar linkage states: If the sum of the shortest and longest link of a planar quadrilateral linkage is less than or equal to the sum of the remaining two links, then the shortest link can rotate fully with respect to a neighboring link. In other words, the condition is satisfied if S + L ≤ P + Q, where S is the shortest link, L is the longest, and P and Q are the other links.
Classification
The movement of a quadrilateral linkage can be classified into eight cases based on the dimensions of its four links. Let a, b, g and h denote the lengths of the input crank, the output crank, the ground link and floating link, respectively. Then, we can construct the three terms:
;
;
.
The movement of a quadrilateral linkage can be classified into eight types based on the positive and negative values for these three terms, T1, T2, and T3.
The cases of T1 = 0, T2 = 0, and T3 = 0 are interesting because the linkages fold. If we distinguish folding quadrilateral linkage, then there are 27 different cases.
The figure shows examples of the various cases for a planar quadrilateral linkage.
The configuration of a quadrilateral linkage may be classified into three types: convex, concave, and crossing. In the convex and concave cases no two links cross over each other. In the crossing linkage two links cross over each other. In the convex case all four internal angles are less than 180 degrees, and in the concave configuration one internal angle is greater than 180 degrees. There exists a simple geometrical relationship between the lengths of the two diagonals of the quadrilateral. For convex and crossing linkages, the length of one diagonal increases if and only if the other decreases. On the other hand, for nonconvex non-crossing linkages, the opposite is the case; one diagonal increases if and only if the other also increases.
Design of four-bar mechanisms
The synthesis, or design, of four-bar mechanisms is important when aiming to produce a desired output motion for a specific input motion. In order to minimize cost and maximize efficiency, a designer will choose the simplest mechanism possible to accomplish the desired motion. When selecting a mechanism type to be designed, link lengths must be determined by a process called dimensional synthesis. Dimensional synthesis involves an iterate-and-analyze methodology which in certain circumstances can be an inefficient process; however, in unique scenarios, exact and detailed procedures to design an accurate mechanism may not exist.
Time ratio
The time ratio (Q) of a four-bar mechanism is a measure of its quick return and is defined as follows:
With four-bar mechanisms there are two strokes, the forward and return, which when added together create a cycle. Each stroke may be identical or have different average speeds. The time ratio numerically defines how fast the forward stroke is compared to the quicker return stroke. The total cycle time () for a mechanism is:
Most four-bar mechanisms are driven by a rotational actuator, or crank, that requires a specific constant speed. This required speed (ωcrank)is related to the cycle time as follows:
Some mechanisms that produce reciprocating, or repeating, motion are designed to produce symmetrical motion. That is, the forward stroke of the machine moves at the same pace as the return stroke. These mechanisms, which are often referred to as in-line design, usually do work in both directions, as they exert the same force in both directions.
Examples of symmetrical motion mechanisms include:
Windshield wipers
Engine mechanisms or pistons
Automobile window crank
Other applications require that the mechanism-to-be-designed has a faster average speed in one direction than the other. This category of mechanism is most desired for design when work is only required to operate in one direction. The speed at which this one stroke operates is also very important in certain machine applications. In general, the return and work-non-intensive stroke should be accomplished as fast as possible. This is so the majority of time in each cycle is allotted for the work-intensive stroke. These quick-return mechanisms are often referred to as offset.
Examples of offset mechanisms include:
Cutting machines
Package-moving devices
With offset mechanisms, it is very important to understand how and to what degree the offset affects the time ratio. To relate the geometry of a specific linkage to the timing of the stroke, an imbalance angle (β) is used. This angle is related to the time ratio, Q, as follows:
Through simple algebraic rearrangement, this equation can be rewritten to solve for β:
Timing charts
Timing charts are often used to synchronize the motion between two or more mechanisms. They graphically display information showing where and when each mechanism is stationary or performing its forward and return strokes. Timing charts allow designers to qualitatively describe the required kinematic behavior of a mechanism.
These charts are also used to estimate the velocities and accelerations of certain four-bar links. The velocity of a link is the time rate at which its position is changing, while the link's acceleration is the time rate at which its velocity is changing. Both velocity and acceleration are vector quantities, in that they have both magnitude and direction; however, only their magnitudes are used in timing charts. When used with two mechanisms, timing charts assume constant acceleration. This assumption produces polynomial equations for velocity as a function of time. Constant acceleration allows for the velocity vs. time graph to appear as straight lines, thus designating a relationship between displacement (ΔR), maximum velocity (vpeak), acceleration (a), and time(Δt). The following equations show this.
Given the displacement and time, both the maximum velocity and acceleration of each mechanism in a given pair can be calculated.
Slider-crank linkage
A slider-crank linkage is a four-bar linkage with three revolute joints and one prismatic, or sliding, joint. The rotation of the crank drives the linear movement the slider, or the expansion of gases against a sliding piston in a cylinder can drive the rotation of the crank.
There are two types of slider-cranks: in-line and offset.
In-line An in-line slider-crank has its slider positioned so the line of travel of the hinged joint of the slider passes through the base joint of the crank. This creates a symmetric slider movement back and forth as the crank rotates.
Offset If the line of travel of the hinged joint of the slider does not pass through the base pivot of the crank, the slider movement is not symmetric. It moves faster in one direction than the other. This is called a quick-return mechanism.
Spherical and spatial four-bar linkages
If the linkage has four hinged joints with axes angled to intersect in a single point, then the links move on concentric spheres and the assembly is called a spherical four-bar linkage. The input-output equations of a spherical four-bar linkage can be applied to spatial four-bar linkages when the variables are replaced by dual numbers. Note that the cited conference paper incorrectly conflates Moore-Penrose pseudoinverses with one-sided inverses of matrices, falsely claiming that the latter are unique whenever they exist. This is contradicted by the fact that admits the set of matrices as all its left inverses.
Bennett's linkage is a spatial four-bar linkage with hinged joints that have their axes angled in a particular way that makes the system movable.
Examples
Other Linkages and Mechanisms
Chebyshev linkage
Chebyshev lambda linkage
Evans "Grasshopper" linkage
Hoecken linkage
Horse-head linkage
Pantograph
Roberts linkage
Valve gear
Watt's linkage
Applications
Bicycle suspension
Biological linkages
Double-beam drawbridge
Double wishbone suspension
Door closer
Foldable steps and foldable chairs
Foot operated machines (grindstone, lathe, sewing machine, treadle, etc.)
Gear shifter
Glider (furniture)
Oscillating fan
Pumpjack
Step-on trash can
Windshield wiper
Simulations
| Technology | Mechanisms | null |
1443028 | https://en.wikipedia.org/wiki/Polar%20vortex | Polar vortex | A circumpolar vortex, or simply polar vortex, is a large region of cold, rotating air; polar vortices encircle both of Earth's polar regions. Polar vortices also exist on other rotating, low-obliquity planetary bodies. The term polar vortex can be used to describe two distinct phenomena; the stratospheric polar vortex, and the tropospheric polar vortex. The stratospheric and tropospheric polar vortices both rotate in the direction of the Earth's spin, but they are distinct phenomena that have different sizes, structures, seasonal cycles, and impacts on weather.
The stratospheric polar vortex is an area of high-speed, cyclonically rotating winds around 15 km to 50 km high, poleward of 50°, and is strongest in winter. It forms during autumn when Arctic or Antarctic temperatures cool rapidly as the polar night begins. The increased temperature difference between the pole and the tropics causes strong winds, and the Coriolis effect causes the vortex to spin up. The stratospheric polar vortex breaks down during spring as the polar night ends. A sudden stratospheric warming (SSW) is an event that occurs when the stratospheric vortex breaks down during winter, and can have significant impacts on surface weather.
The tropospheric polar vortex is often defined as the area poleward of the tropospheric jet stream. The equatorward edge is around 40° to 50°, and it extends from the surface up to around 10 km to 15 km. Its yearly cycle differs from the stratospheric vortex because the tropospheric vortex exists all year, but is similar to the stratospheric vortex since it is also strongest in winter when the polar regions are coldest.
The tropospheric polar vortex was first described as early as 1853. The stratospheric vortex's SSWs were discovered in 1952 with radiosonde observations at altitudes higher than 20 km. The tropospheric polar vortex was mentioned frequently in the news and weather media in the cold North American winter of 2013–2014, popularizing the term as an explanation of very cold temperatures. The tropospheric vortex increased in public visibility in 2021 as a result of extreme frigid temperatures in the central United States, with newspapers linking its effects to climate change.
Ozone depletion occurs most heavily within the polar vortices – particularly over the Southern Hemisphere – reaching a maximum depletion in the spring.
Snow or ice weather phenomena
Arctic and Antarctic vortices
Northern Hemisphere
When the tropospheric vortex of the Arctic is strong, it has a well defined and nearly circular shape. There is a single vortex with a jet stream that is well constrained near the polar front, and the Arctic air is well contained. When this northern tropospheric vortex weakens, it breaks into two or more smaller vortices, the strongest of which are near Baffin Island, Nunavut, and the others over northeast Siberia. When it is very weak, the flow of Arctic air becomes more disorganized, and masses of cold Arctic air can push equatorward, bringing with them a rapid and sharp temperature drop.
A deep freeze that gripped much of the United States and Canada in late January 2019 has been blamed on a "polar vortex". This is not the scientifically correct use of the term polar vortex, but instead is referring to outbreaks of cold Arctic air caused by a weakened polar vortex. The US National Weather Service warned that frostbite is possible within just 10 minutes of being outside in such extreme temperatures, and hundreds of schools, colleges, and universities in the affected areas were closed. Around 21 people died in US due to severe frostbite. States within the midwest region of the United States had windchills just above -50 °F (-45 °C). The Polar vortex is also thought to have had effects in Europe. For example, the 2013–14 United Kingdom winter floods were blamed on the Polar vortex bringing severe cold in the United States and Canada. Similarly, the severe cold in the United Kingdom in the winters of 2009–10 and 2010–11 were also blamed on the Polar vortex.
Southern Hemisphere
The Antarctic vortex of the Southern Hemisphere is a single low-pressure zone that is found near the edge of the Ross ice shelf, near 160 west longitude. When the polar vortex is strong, the mid-latitude Westerlies (winds at the surface level between 30° and 60° latitude from the west) increase in strength and are persistent. When the polar vortex is weak, high-pressure zones of the mid-latitudes may push poleward, moving the polar vortex, jet stream, and polar front equatorward. The jet stream is seen to "buckle" and deviate south. This rapidly brings cold dry air into contact with the warm, moist air of the mid-latitudes, resulting in a rapid and dramatic change of weather known as a "cold snap".
In Australia, the polar vortex, known there as a "polar blast" or "polar plunge", is a cold front that drags air from Antarctica which brings rain showers, snow (typically inland, with blizzards occurring in the highlands), gusty icy winds, and hail in the south-eastern parts of the country, such as in Victoria, Tasmania, the southeast coast of South Australia and the southern half of New South Wales (but only on the windward side of the Great Dividing Range, whereas the leeward side will be affected by foehn winds).
Identification
The bases of the two polar vortices are located in the middle and upper troposphere and extend into the stratosphere. Beneath that lies a large mass of cold, dense Arctic air. The interface between the cold dry air mass of the pole and the warm moist air mass farther south defines the location of the polar front. The polar front is centered roughly at 60° latitude. A polar vortex strengthens in the winter and weakens in the summer because of its dependence on the temperature difference between the equator and the poles.
Polar cyclones are low-pressure zones embedded within the polar air masses, and exist year-round. The stratospheric polar vortex develops at latitudes above the subtropical jet stream. Horizontally, most polar vortices have a radius of less than . Since polar vortices exist from the stratosphere downward into the mid-troposphere, a variety of heights/pressure levels are used to mark its position. The 50 hPa pressure surface is most often used to identify its stratospheric location. At the level of the tropopause, the extent of closed contours of potential temperature can be used to determine its strength. Others have used levels down to the 500 hPa pressure level (about above sea level during the winter) to identify the polar vortex.
Duration and strength
Polar vortices are weakest during summer and strongest during winter. Extratropical cyclones that migrate into higher latitudes when the polar vortex is weak can disrupt the single vortex creating smaller vortices (cold-core lows) within the polar air mass. Those individual vortices can persist for more than a month.
Volcanic eruptions in the tropics can lead to a stronger polar vortex during winter for as long as two years afterwards. The strength and position of the polar vortex shapes the flow pattern in a broad area about it. An index which is used in the northern hemisphere to gauge its magnitude is the Arctic oscillation.
When the Arctic vortex is at its strongest, there is a single vortex, but normally, the Arctic vortex is elongated in shape, with two cyclone centers, one over Baffin Island in Canada and the other over northeast Siberia. When the Arctic pattern is at its weakest, subtropic air masses can intrude poleward causing the Arctic air masses to move equatorward, as during the Winter 1985 Arctic outbreak. The Antarctic polar vortex is more pronounced and persistent than the Arctic one. In the Arctic the distribution of land masses at high latitudes in the Northern Hemisphere gives rise to Rossby waves which contribute to the breakdown of the polar vortex, whereas in the Southern Hemisphere the vortex is less disturbed. The breakdown of the polar vortex is an extreme event known as a sudden stratospheric warming, here the vortex completely breaks down and an associated warming of 30–50 °C (54–90 °F) over a few days can occur.
The waxing and waning of the polar vortex is driven by the movement of mass and the transfer of heat in the polar region. In the autumn, the circumpolar winds increase in speed and the polar vortex rises into the stratosphere. The result is that the polar air forms a coherent rotating air mass: the polar vortex. As winter approaches, the vortex core cools, the winds decrease, and the vortex energy declines. Once late winter and early spring approach the vortex is at its weakest. As a result, during late winter, large fragments of the vortex air can be diverted into lower latitudes by stronger weather systems intruding from those latitudes. In the lowest level of the stratosphere, strong potential vorticity gradients remain, and the majority of that air remains confined within the polar air mass into December in the Southern Hemisphere and April in the Northern Hemisphere, well after the breakup of the vortex in the mid-stratosphere.
The breakup of the northern polar vortex occurs between mid March to mid May. This event signifies the transition from winter to spring, and has impacts on the hydrological cycle, growing seasons of vegetation, and overall ecosystem productivity. The timing of the transition also influences changes in sea ice, ozone, air temperature, and cloudiness. Early and late polar breakup episodes have occurred, due to variations in the stratospheric flow structure and upward spreading of planetary waves from the troposphere. As a result of increased waves into the vortex, the vortex experiences more rapid warming than normal, resulting in an earlier breakup and spring. When the breakup comes early, it is characterized by with persistent of remnants of the vortex. When the breakup is late, the remnants dissipate rapidly. When the breakup is early, there is one warming period from late February to middle March. When the breakup is late, there are two warming periods, one January, and one in March. Zonal mean temperature, wind, and geopotential height exert varying deviations from their normal values before and after early breakups, while the deviations remain constant before and after late breakups. Scientists are connecting a delay in the Arctic vortex breakup with a reduction of planetary wave activities, few stratospheric sudden warming events, and depletion of ozone.
Sudden stratospheric warming events are associated with weaker polar vortices. This warming of stratospheric air can reverse the circulation in the Arctic Polar Vortex from counter-clockwise to clockwise. These changes aloft force changes in the troposphere below. An example of an effect on the troposphere is the change in speed of the Atlantic Ocean circulation pattern. A soft spot just south of Greenland is where the initial step of downwelling occurs, nicknamed the "Achilles Heel of the North Atlantic". Small amounts of heating or cooling traveling from the polar vortex can trigger or delay downwelling, altering the Gulf Stream Current of the Atlantic, and the speed of other ocean currents. Since all other oceans depend on the Atlantic Ocean's movement of heat energy, climates across the planet can be dramatically affected. The weakening or strengthening of the polar vortex can alter the sea circulation more than a mile beneath the waves. Strengthening storm systems within the troposphere that cool the poles, intensify the polar vortex. La Niña–related climate anomalies significantly strengthen the polar vortex. Intensification of the polar vortex produces changes in relative humidity as downward intrusions of dry, stratospheric air enter the vortex core. With a strengthening of the vortex comes a longwave cooling due to a decrease in water vapor concentration near the vortex. The decreased water content is a result of a lower tropopause within the vortex, which places dry stratospheric air above moist tropospheric air. Instability is caused when the vortex tube, the line of concentrated vorticity, is displaced. When this occurs, the vortex rings become more unstable and prone to shifting by planetary waves. The planetary wave activity in both hemispheres varies year-to-year, producing a corresponding response in the strength and temperature of the polar vortex. The number of waves around the perimeter of the vortex are related to the core size; as the vortex core decreases, the number of waves increase.
The degree of the mixing of polar and mid-latitude air depends on the evolution and position of the polar night jet. In general, the mixing is less inside the vortex than outside. Mixing occurs with unstable planetary waves that are characteristic of the middle and upper stratosphere in winter. Prior to vortex breakdown, there is little transport of air out of the Arctic Polar Vortex due to strong barriers above 420 km (261 miles). The polar night jet which exists below this, is weak in the early winter. As a result, it does not deviate any descending polar air, which then mixes with air in the mid-latitudes. In the late winter, air parcels do not descend as much, reducing mixing. After the vortex is broken up, the ex-vortex air is dispersed into the middle latitudes within a month.
Sometimes, a mass of the polar vortex breaks off before the end of the final warming period. If large enough, the piece can move into Canada and the Midwestern, Central, Southern, and Northeastern United States. This diversion of the polar vortex can occur due to the displacement of the polar jet stream; for example, the significant northwestward direction of the polar jet stream in the western part of the United States during the winters of 2013–2014, and 2014–2015. This caused warm, dry conditions in the west, and cold, snowy conditions in the north-central and northeast. Occasionally, the high-pressure air mass, called the Greenland Block, can cause the polar vortex to divert to the south, rather than follow its normal path over the North Atlantic.
Extreme weather
A study in 2001 found that stratospheric circulation can have anomalous effects on weather regimes. In the same year, researchers found a statistical correlation between weak polar vortex and outbreaks of severe cold in the Northern Hemisphere. In later years, scientists identified interactions with Arctic sea ice decline, reduced snow cover, evapotranspiration patterns, NAO anomalies or weather anomalies which are linked to the polar vortex and jet stream configuration.
Climate change
Ozone depletion
The chemistry of the Antarctic polar vortex has created severe ozone depletion, although the effect has been weakening since the 2000s. It is expected to return to 1980 levels in about 2075. The nitric acid in polar stratospheric clouds reacts with chlorofluorocarbons to form chlorine, which catalyzes the photochemical destruction of ozone. Chlorine concentrations build up during the polar winter, and the consequent ozone destruction is greatest when the sunlight returns in spring. These clouds can only form at temperatures below about .
Since there is greater air exchange between the Arctic and the mid-latitudes, ozone depletion at the north pole is much less severe than at the south. Accordingly, the seasonal reduction of ozone levels over the Arctic is usually characterized as an "ozone dent", whereas the more severe ozone depletion over the Antarctic is considered an "ozone hole". That said, chemical ozone destruction in the 2011 Arctic polar vortex attained, for the first time, a level clearly identifiable as an Arctic "ozone hole".
Outside Earth
Other astronomical bodies are also known to have polar vortices, including Venus (double vortex – that is, two polar vortices at a pole), Mars, Jupiter, Saturn, and Saturn's moon Titan.
Saturn's south pole is the only known hot polar vortex in the solar system.
| Physical sciences | Atmospheric circulation | null |
29113700 | https://en.wikipedia.org/wiki/Autism | Autism | Autism spectrum disorder (ASD), or simply autism, is a neurodevelopmental disorder "characterized by persistent deficits in social communication and social interaction across multiple contexts" and "restricted, repetitive patterns of behavior, interests, or activities". Sensory abnormalities are also included in the diagnostic manuals. Common associated traits such as motor coordination impairment are typical of the condition but not required for diagnosis. A formal diagnosis requires that symptoms cause significant impairment in multiple functional domains; in addition, the symptoms must be atypical or excessive for the person's age and sociocultural context.
Autism is a spectrum, meaning it manifests in various ways, with its severity and support needs varying widely across different autistic people. For example, some autistic people are nonverbal, while others have proficient spoken language. Furthermore, the spectrum is multi-dimensional and not all dimensions have been identified .
Public health authorities and guideline developers classify autism as a neurodevelopmental disorder, but the autism rights movement (and some researchers) disagree with the classification. From the latter point of view, autistic people may be diagnosed with a disability, but that disability may be rooted in the structures of a society rather than the person. On the contrary, other scientists argue that autism impairs functioning in many ways that are inherent to the disorder itself and unrelated to society. The neurodiversity perspective has led to significant controversy among those who are autistic and advocates, practitioners, and charities.
The precise causes of autism are unknown in most individual cases. Research shows that the disorder is highly heritable and polygenic, and neurobiological risks from the environment are also relevant. Boys are also significantly far more frequently diagnosed than girls.
Autism frequently co-occurs with attention deficit hyperactivity disorder (ADHD), epilepsy, and intellectual disability. Disagreements persist about what should be part of the diagnosis, whether there are meaningful subtypes or stages of autism, and the significance of autism-associated traits in the wider population.
The combination of broader criteria, increased awareness, and the potential increase of actual prevalence has led to considerably increased estimates of autism prevalence since the 1990s. The WHO estimates about 1 in 100 children had autism between 2012 and 2021, as that was the average estimate in studies during that period, with a trend of increasing prevalence over time. This increasing prevalence has contributed to the myth perpetuated by anti-vaccine activists that autism is caused by vaccines.
There is no known cure for autism. Some advocates dispute the need to find one. Interventions such as applied behavior analysis (ABA), speech therapy, and occupational therapy can help these children gain self-care, social, and language skills. Guidelines from the Centres for Disease Control (CDC), and European Society for Child & Adolescent Psychiatry endorse the use of ABA on the grounds that it reduces symptoms impairing daily functioning and quality of life, but the National Institute for Health and Care Excellence cites a lack of high-quality evidence to support its use. Additionally, some in the autism rights movement oppose its application due to a perception that it emphasises normalisation. No medication has been shown to reduce ASD's core symptoms, but some can alleviate comorbid issues.
Classification
Spectrum model
Before the DSM-5 (2013) and ICD-11 CDDR (2024) diagnostic manuals were adopted, ASD was found under the diagnostic category pervasive developmental disorder. The previous system relied on a set of closely related and overlapping diagnoses such as Asperger syndrome and the syndrome formerly known as Kanner syndrome. This created unclear boundaries between the terms, so for the DSM-5 and CDDR, a spectrum approach was taken. The new system is also more restrictive, meaning fewer people qualify for diagnosis.
The DSM-5 and CDDR use different categorization tools to define this spectrum. DSM-5 uses a "level" system, which ranks how in need of support the patient is, level 1 being the mildest and level 3 the severest, while the CDDR system has two axes, intellectual impairment and language impairment, as these are seen as the most crucial factors.
Autism is currently defined as a highly variable neurodevelopmental disorder that is generally thought to cover a broad and deep spectrum, manifesting very differently from one person to another. Some have high support needs, may be nonspeaking, and experience developmental delays; this is more likely with other co-existing diagnoses. Others have relatively low support needs; they may have more typical speech-language and intellectual skills but atypical social/conversation skills, narrowly focused interests, and wordy, pedantic communication. They may still require significant support in some areas of their lives. The spectrum model should not be understood as a continuum running from mild to severe, but instead means that autism can present very differently in each person. How it presents in a person can depend on context, and may vary over time.
While the DSM and ICD greatly influence each other, there are also differences. For example, Rett syndrome was included in ASD in the DSM-5, but in the ICD-11 it was excluded and placed in the chapter on Developmental Anomalies. The ICD and the DSM change over time, and there has been collaborative work toward a convergence of the two since 1980 (when DSM-III was published and ICD-9 was current), including more rigorous biological assessment—in place of historical experience—and a simplification of the classification system.
As of 2023, empirical and theoretical research is leading to a growing consensus among researchers that the established ASD criteria are ineffective descriptors of autism as a whole, and that alternative research approaches must be encouraged, such as going back to autism prototypes, exploring new causal models of autism, or developing transdiagnostic endophenotypes. Proposed alternatives to the current disorder-focused spectrum model deconstruct autism into at least two separate phenomena: (1) a non-pathological spectrum of behavioral traits in the population, and (2) the neuropathological burden of rare genetic mutations and environmental risk factors potentially leading to neurodevelopmental and psychological disorders, (3) governed by an individual's cognitive ability to compensate.
ICD
The World Health Organization's International Classification of Diseases (11th revision), ICD-11, was released in June 2018 and came into full effect as of January 2022. It describes ASD as follows:
ICD-11 was produced by professionals from 55 countries out of the 90 involved and is the most widely used reference worldwide.
DSM
The American Psychiatric Association's Diagnostic and Statistical Manual of Mental Disorders, Fifth Edition, Text Revision (DSM-5-TR), released in 2022, is the current version of the DSM. It is the predominant mental health diagnostic system used in the United States and Canada, and is often used in Anglophone countries.
Its fifth edition, DSM-5, released in May 2013, was the first to define ASD as a single diagnosis, which is still the case in the DSM-5-TR. ASD encompasses previous diagnoses, including the four traditional diagnoses of autism—classic autism, Asperger syndrome, childhood disintegrative disorder, and pervasive developmental disorder not otherwise specified (PDD-NOS)—and the range of diagnoses that included the word "autism". Rather than distinguishing among these diagnoses, the DSM-5 and DSM-5-TR adopt a dimensional approach with one diagnostic category for disorders that fall under the autism spectrum umbrella. Within that category, the DSM-5 and the DSM include a framework that differentiates each person by dimensions of symptom severity, as well as by associated features (i.e., the presence of other disorders or factors that likely contribute to the symptoms, other neurodevelopmental or mental disorders, intellectual disability, or language impairment). The symptom domains are (a) social communication and (b) restricted, repetitive behaviors, and there is the option of specifying a separate severity—the negative effect of the symptoms on the person—for each domain, rather than just overall severity. Before the DSM-5, the DSM separated social deficits and communication deficits into two domains. Further, the DSM-5 changed to an onset age in the early developmental period, with a note that symptoms may manifest later when social demands exceed capabilities, rather than the previous, more restricted three years of age. These changes remain in the DSM-5-TR.
Signs and symptoms
Pre-diagnosis
For many autistic people, characteristics first appear during infancy or childhood and follow a steady course without remission (different developmental timelines are described in more detail below). Autistic people may be severely impaired in some respects but average, or even superior, in others.
Clinicians consider assessment for ASD when a patient shows:
Regular difficulties in social interaction or communication
Restricted or repetitive behaviors (often called "stimming")
Resistance to changes or restricted interests
These features are typically assessed with the following, when appropriate:
Problems in obtaining or sustaining employment or education
Difficulties in initiating or sustaining social relationships
Connections with mental health or learning disability services
A history of neurodevelopmental conditions (including learning disabilities and ADHD) or mental health conditions
There are many signs associated with autism; the presentation varies widely. Common signs and symptoms include:
Abnormalities in eye contact
Little or no babbling as an infant
Not showing interest in indicated objects
Delayed language skills (e.g., having a smaller vocabulary than peers or difficulty expressing themselves in words)
Reduced interest in other children or caretakers, possibly with more interest in objects
Difficulty playing reciprocal games (e.g., peek-a-boo)
Hyper- or hypo-sensitivity to or unusual response to the smell, texture, sound, taste, or appearance of things
Resistance to changes in routine
Repetitive, limited, or otherwise unusual usage of toys (e.g., lining up toys)
Repetition of words or phrases, including echolalia
Repetitive motions or movements, including stimming
Broader autism phenotype
The broader autism phenotype describes people who may not have ASD but do have autistic traits, such as abnormalities in eye contact and stimming.
Social and communication skills
According to the medical model, autistic people experience social communications impairments. Until 2013, deficits in social function and communication were considered two separate symptom domains. The current social communication domain criteria for autism diagnosis require people to have deficits across three social skills: social-emotional reciprocity, nonverbal communication, and developing and sustaining relationships.
Social-emotional reciprocity
A deficit-based view predicts that autistic–autistic interaction would be less effective than autistic–non-autistic interactions or even non-functional. But recent research has found that autistic–autistic interactions are as effective in information transfer as interactions between non-autistics are, and that communication breaks down only between autistics and non-autistics. Also contrary to social cognitive deficit interpretations, recent (2019) research recorded similar social cognitive performances in autistic and non-autistic adults, with both of them rating autistic individuals less favorably than non-autistic individuals; however, autistic individuals showed more interest in engaging with autistic people than non-autistic people did, and learning of a person's ASD diagnosis did not influence their interest level.
Thus, there has been a recent shift to acknowledge that autistic people may simply respond and behave differently than people without ASD. So far, research has identified two unconventional features by which autistic people create shared understanding (intersubjectivity): "a generous assumption of common ground that, when understood, led to rapid rapport, and, when not understood, resulted in potentially disruptive utterances; and a low demand for coordination that ameliorated many challenges associated with disruptive turns." Autistic interests, and thus conversational topics, seem to be largely driven by an intense interest in specific topics (monotropism).
Historically, autistic children were said to be delayed in developing a theory of mind, and the empathizing–systemizing theory has argued that while autistic people have compassion (affective empathy) for others with similar presentation of symptoms, they have limited, though not necessarily absent, cognitive empathy. This may present as social naïvety, lower than average intuitive perception of the utility or meaning of body language, social reciprocity, or social expectations, including the habitus, social cues, and some aspects of sarcasm, which to some degree may also be due to comorbid alexithymia. But recent research has increasingly questioned these findings, as the "double empathy problem" theory (2012) argues that there is a lack of mutual understanding and empathy between both non-autistic persons and autistic individuals.
As communication is bidirectional, research on communication difficulties has since also begun to study non-autistic behavior, with researcher Catherine Crompton writing in 2020 that non-autistic people "struggle to identify autistic mental states, identify autistic facial expressions, overestimate autistic egocentricity, and are less willing to socially interact with autistic people. Thus, although non-autistic people are generally characterised as socially skilled, these skills may not be functional or effectively applied when interacting with autistic people." Any previously observed communication deficits of autistic people may thus have been constructed through a neurotypical bias in autism research, which has come to be scrutinized for "dehumanization, objectification, and stigmatization". Recent research has proposed that autistics' lack of readability and a neurotypical lack of effort to interpret atypical signals may cause a negative interaction loop, increasingly driving both groups apart into two distinct groups with different social interaction styles.
Verbal, minimally verbal, or nonverbal communication
Differences in verbal communication begin to be noticeable in childhood, as many autistic children develop language skills at an uneven pace. Verbal communication may be delayed or never developed (nonverbal autism), while reading ability may be present before school age (hyperlexia). Reduced joint attention seem to distinguish autistic from non-autistic infants. Infants may show delayed onset of babbling, unusual gestures, diminished responsiveness, and vocal patterns that are not synchronized with the caregiver. In their second and third years, autistic children may have less frequent and less diverse babbling, consonants, words, and word combinations, and their gestures may be less often integrated with words. Autistic children are less likely to make requests or share experiences and more likely to simply repeat others' words (echolalia). The CDC estimated in 2015 that around 40% of autistic children do not speak at all. Autistic adults' verbal communication skills largely depend on when and how well speech is acquired during childhood.
Autistic people display atypical nonverbal behaviors or show differences in nonverbal communication. They may make infrequent eye contact, even when called by name, or avoid it altogether. This may be due to the high amount of sensory input received when making eye contact. Autistic people often recognize fewer emotions and their meaning from others' facial expressions, and may not respond with facial expressions expected by their non-autistic peers. Temple Grandin, an autistic woman involved in autism activism, described her inability to understand neurotypicals' social communication as leaving her feeling "like an anthropologist on Mars". Autistic people struggle to understand the social context and subtext of neurotypical conversational or printed situations, and form different conclusions about the content. Autistic people may not control the volume of their voice in different social settings. At least half of autistic children have atypical prosody.
Developing and sustaining relationships
What may look like self-involvement or indifference to non-autistic people stems from autistic differences in recognizing how other people have their own personalities, perspectives, and interests. Most published research focuses on the interpersonal relationship difficulties between autistic people and their non-autistic counterparts and how to solve them through teaching neurotypical social skills, but newer research has also evaluated what autistic people want from friendships, such as a sense of belonging and good mental health. Children with ASD are more frequently involved in bullying situations than their non-autistic peers, and predominantly experience bullying as victims rather than perpetrators or victim-perpetrators, especially after controlling for comorbid psychopathology. Prioritizing dependability and intimacy in friendships during adolescence, coupled with lowered friendship quantity and quality, often lead to increased loneliness in autistic people. As they progress through life, autistic people observe and form a model of social patterns, and develop coping mechanisms, referred to as "masking", which have recently been found to come with psychological costs and a higher increased risk of suicidality.
Restricted and repetitive behaviors
ASD includes a wide variety of characteristics. Some of these include behavioral characteristics, which widely range from slow development of social and learning skills to difficulties creating connections with other people. Autistic people may experience these challenges with forming connections due to anxiety or depression, which they are more likely to experience, and may respond by isolating themselves.
Other behavioral characteristics include abnormal responses to sensations (such as sights, sounds, touch, taste and smell) and problems keeping a consistent speech rhythm. The latter problem influences social skills, leading to potential problems in understanding for interlocutors. Autistic people's behavioral characteristics typically influence development, language, and social competence. Their behavioral characteristics can be observed as perceptual disturbances, disturbances of development rate, relating, speech and language, and motility.
The second core symptom of autism spectrum is a pattern of restricted and repetitive behaviors, activities, and interests. In order to be diagnosed with ASD under the DSM-5-TR, a person must have at least two of the following behaviors:
Repetitive behaviors – Repetitive behaviors such as rocking, hand flapping, finger flicking, head banging, or repeating phrases or sounds. These behaviors may occur constantly or only when the person gets stressed, anxious, or upset. These behaviors are also known as stimming.
Resistance to change – A strict adherence to routines such as eating certain foods in a specific order or taking the same path to school every day. The person may become distressed if there is a change or disruption to their routine.
Restricted interests – An excessive interest in a particular activity, topic, or hobby, and devoting all their attention to it. For example, young children might completely focus on things that spin and ignore everything else. Older children might try to learn everything about a single topic, such as the weather or sports, and perseverate or talk about it constantly.
Sensory reactivity – An unusual reaction to certain sensory inputs, such as negative reaction to specific sounds or textures, fascination with lights or movements, or apparent indifference to pain or heat.
Autistic people can display many forms of repetitive or restricted behavior, which the Repetitive Behavior Scale-Revised (RBS-R) categorizes as follows.
Stereotyped behaviors: Repetitive movements, such as hand flapping, head rolling, or body rocking.
Compulsive behaviors: Time-consuming behaviors intended to reduce anxiety, that a person feels compelled to perform repeatedly or according to rigid rules, such as placing objects in a specific order, checking things, or handwashing.
Sameness: Resistance to change; for example, insisting that the furniture not be moved or refusing to be interrupted.
Ritualistic behavior: Unvarying pattern of daily activities, such as an unchanging menu or a dressing ritual. This is closely associated with sameness and an independent validation has suggested combining the two factors.
Self-injurious behaviors: Behaviors such as eye-poking, skin-picking, hand-biting and head-banging.
Mental health, self-injury and suicide
Self-injurious behaviors are relatively common in autistic people, and can include head-banging, self-cutting, self-biting, and hair-pulling. Some of these can result in serious injury or death. Autistic people are about three times as likely as non-autistic people to engage in self-injury.
Theories about the cause of self-injurious behavior in children with developmental delay, including autistic children, include:
Frequency or continuation of self-injurious behavior can be influenced by environmental factors (e.g., reward in return for halting self-injurious behavior). This theory does not apply to younger children with autism. There is some evidence that frequency of self-injurious behavior can be reduced by removing or modifying environmental factors that reinforce the behavior.
Higher rates of self-injury are noted in socially isolated autistic people. Studies have shown that socialization skills are related factors to self-injurious behavior for autistic people.
Self-injury could be a response to modulate pain perception when chronic pain or other health problems that cause pain are present.
Abnormal basal ganglia connectivity may predispose to self-injurious behavior.
Risk factors for self-harm and suicidality include circumstances that could affect anyone, such as mental health problems (e.g., anxiety disorder) and social problems (e.g., unemployment and social isolation), plus factors that affect only autistic people, such as actively trying to behave like a neurotypical person, which is called masking. Approximately 8 in 10 people with autism suffer from a mental health problem in their lifetime, in comparison to 1 in 4 of the general population that suffers from a mental health problem in their lifetimes. A 2019 meta-analysis identified autistic people as being four times more likely to have depression than non-autistic people, with approximately 40% of autistic adults having depression.
Rates of suicidality vary significantly depending upon what is being measured. This is partly because questionnaires developed for neurotypical subjects are not always valid for autistic people. As of 2023, the Suicidal Behaviours Questionnaire–Autism Spectrum Conditions (SBQ-ASC) is the only test validated for autistic people. According to some estimates, about a quarter of autistic youth and a third of all autistic people have experienced suicidal ideation at some point. Rates of suicidal ideation are the same for people formally diagnosed with autism and people who have typical intelligence and are believed to have autism but have not been diagnosed. The suicide rate for verbal autistics is nine times that of the general population. In a 2014 study of late diagnosed autistic adults, 66% had experienced suicidal ideation (nine times higher than the general population) and 35% had a suicide plan or had made a suicide attempt.
Although most people who attempt suicide are not autistic, autistic people are about three times as likely as non-autistic people to make a suicide attempt. Less than 10% of autistic youth have attempted suicide, but 15% to 25% autistic adults have. The rates of suicide attempts are the same among people formally diagnosed with autism and those who have typical intelligence and are believed to have autism but have not been diagnosed. The suicide risk is lower among cisgender autistic males and autistic people with intellectual disabilities. The rate of suicide results in a global excess mortality among autistic people equal to approximately 2% of all suicide deaths each year.
Burnout
Studies have supported the common belief that autistic people become exhausted or burnt out in some situations.
Other features
Autistic people may have symptoms that do not contribute to the official diagnosis, but that can affect the person or the family.
Some autistic people show unusual or notable abilities, ranging from splinter skills (such as the memorization of trivia) to rare talents in mathematics, music, or artistic reproduction, which in exceptional cases are considered a part of the savant syndrome. One study describes how some autistic people show superior skills in perception and attention relative to the general population. Sensory abnormalities are found in over 90% of autistic people, and are considered core features by some.
More generally, autistic people tend to show a "spiky skills profile", with strong abilities in some areas contrasting with much weaker abilities in others.
Autistic people are less likely to show cognitive or emotional biases, and usually process information more rationally. On the other hand, most autistic people exhibit lower levels of emotional intelligence, the ability to understand nonverbal clues about other people's feelings.
Differences between the previously recognized disorders under the autism spectrum are greater for under-responsivity (for example, walking into things) than for over-responsivity (for example, distress from loud noises) or for sensation seeking (for example, rhythmic movements). An estimated 60–80% of autistic people have motor signs that include poor muscle tone, poor motor planning, and toe walking; deficits in motor coordination are pervasive across ASD and are greater in autism proper.
Pathological demand avoidance can occur. People with this set of autistic symptoms are more likely to refuse to do what is asked or expected of them, even to activities they enjoy.
Unusual or atypical eating behavior occurs in about three-quarters of children with ASD, to the extent that it was formerly a diagnostic indicator. Selectivity is the most common problem, although eating rituals and food refusal also occur.
Problematic digital media use
Possible causes
Exactly what causes autism remains unknown. It was long mostly presumed that there is a common cause at the genetic, cognitive, and neural levels for the social and non-social components of ASD's symptoms, described as a triad in the classic autism criteria. But it is increasingly suspected that autism is instead a complex disorder whose core aspects have distinct causes that often cooccur. It is unlikely that ASD has a single cause; many risk factors identified in the research literature may contribute to ASD. These include genetics, prenatal and perinatal factors (meaning factors during pregnancy or very early infancy), neuroanatomical abnormalities, and environmental factors. It is possible to identify general factors, but much more difficult to pinpoint specific ones. Given the current state of knowledge, prediction can only be of a global nature and so requires the use of general markers.
Biological subgroups
Research into causes has been hampered by the inability to identify biologically meaningful subgroups within the autistic population and by the traditional boundaries between the disciplines of psychiatry, psychology, neurology and pediatrics. Newer technologies such as fMRI and diffusion tensor imaging can help identify biologically relevant phenotypes (observable traits) that can be viewed on brain scans, to help further neurogenetic studies of autism; one example is lowered activity in the fusiform face area of the brain, which is associated with impaired perception of people versus objects. It has been proposed to classify autism using genetics as well as behavior.
Syndromic autism and non-syndromic autism
Autism spectrum disorder (ASD) can be classified into two categories: "syndromic autism" and "non-syndromic autism".
Syndromic autism refers to cases where ASD is one of the characteristics associated with a broader medical condition or syndrome, representing about 25% of ASD cases. The causes of syndromic autism are often known, and monogenic disorders account for approximately 5% of these cases.
Non-syndromic autism, also known as classic or idiopathic autism, represents the majority of cases, and its cause is typically polygenic and unknown.
Genetics
Autism has a strong genetic basis, although the genetics of autism are complex and it is unclear whether ASD is explained more by rare mutations with major effects, or by rare multi-gene interactions of common genetic variants. Complexity arises due to interactions among multiple genes, the environment, and epigenetic factors which do not change DNA sequencing but are heritable and influence gene expression. Many genes have been associated with autism through sequencing the genomes of affected people and their parents. But most of the mutations that increase autism risk have not been identified. Typically, autism cannot be traced to a Mendelian (single-gene) mutation or to a single chromosome abnormality, and none of the genetic syndromes associated with ASD have been shown to selectively cause ASD. Numerous genes have been found, with only small effects attributable to any particular gene. Most loci individually explain less than 1% of cases of autism. , it appeared that between 74% and 93% of ASD risk is heritable. After an older child is diagnosed with ASD, 7% to 20% of subsequent children are likely to be as well. If parents have one autistic child, they have a 2% to 8% chance of having a second child who is autistic. If the autistic child is an identical twin, the other will be affected 36% to 95% of the time. A fraternal twin is affected up to 31% of the time. The large number of autistic people with unaffected family members may result from spontaneous structural variation, such as deletions, duplications or inversions in genetic material during meiosis. Hence, a substantial fraction of autism cases may be traceable to genetic causes that are highly heritable but not inherited: that is, the mutation that causes the autism is not present in the parental genome.
, understanding of genetic risk factors had shifted from a focus on a few alleles to an understanding that genetic involvement in ASD is probably diffuse, depending on a large number of variants, some of which are common and have a small effect, and some of which are rare and have a large effect. The most common gene disrupted with large effect rare variants appeared to be CHD8, but less than 0.5% of autistic people have such a mutation. The gene CHD8 encodes the protein chromodomain helicase DNA binding protein 8, which is a chromatin regulator enzyme that is essential during fetal development. CHD8 is an adenosine triphosphate (ATP)–dependent enzyme. The protein contains an Snf2 helicase domain that is responsible for the hydrolysis of ATP to adenosine diphosphate (ADP). CHD8 encodes a DNA helicase that functions as a repressor of transcription, remodeling chromatin structure by altering the position of nucleosomes. CHD8 negatively regulates Wnt signaling. Wnt signaling is important in the vertebrate early development and morphogenesis. It is believed that CHD8 also recruits the linker histone H1 and causes the repression of β-catenin and p53 target genes. The importance of CHD8 can be observed in studies where CHD8-knockout mice died after 5.5 embryonic days because of widespread p53-induced apoptosis. Some studies have determined the role of CHD8 in autism spectrum disorder (ASD). CHD8 expression significantly increases during human mid-fetal development. The chromatin remodeling activity and its interaction with transcriptional regulators have shown to play an important role in ASD aetiology. The developing mammalian brain has conserved CHD8 target regions that are associated with ASD risk genes. The knockdown of CHD8 in human neural stem cells results in dysregulation of ASD risk genes that are targeted by CHD8. Recently CHD8 has been associated with the regulation of long non-coding RNAs (lncRNAs), and the regulation of X chromosome inactivation (XCI) initiation, via regulation of Xist long non-coding RNA, the master regulator of XCI, though competitive binding to Xist regulatory regions.
Some ASD is associated with clearly genetic conditions, like fragile X syndrome, but only around 2% of autistic people have fragile X. Hypotheses from evolutionary psychiatry suggest that these genes persist because they are linked to human inventiveness, intelligence or systemising.
Current research suggests that genes that increase susceptibility to ASD are ones that control protein synthesis in neuronal cells in response to cell needs, activity and adhesion of neuronal cells, synapse formation and remodeling, and excitatory to inhibitory neurotransmitter balance. Therefore, although up to 1,000 different genes are thought to increase the risk of ASD, all of them eventually affect normal neural development and connectivity between different functional areas of the brain in a similar manner that is characteristic of an ASD brain. Some of these genes are known to modulate production of the GABA neurotransmitter, the nervous system's main inhibitory neurotransmitter. These GABA-related genes are under-expressed in an ASD brain. On the other hand, genes controlling expression of glial and immune cells in the brain, e.g. astrocytes and microglia, respectively, are overexpressed, which correlates with increased number of glial and immune cells found in postmortem ASD brains. Some genes under investigation in ASD pathophysiology are those that affect the mTOR signaling pathway, which supports cell growth and survival.
All these genetic variants contribute to the development of the autism spectrum, but it cannot be guaranteed that they are determinants for the development.
ASD may be under-diagnosed in women and girls due to an assumption that it is primarily a male condition, but genetic phenomena such as imprinting and X linkage have the ability to raise the frequency and severity of conditions in males, and theories have been put forward for a genetic reason why males are diagnosed more often, such as the imprinted brain hypothesis and the extreme male brain theory.
Early life
Several prenatal and perinatal complications have been reported as possible risk factors for autism. These risk factors include maternal gestational diabetes, maternal and paternal age over 30, bleeding during pregnancy after the first trimester, use of certain prescription medication (e.g. valproate) during pregnancy, and meconium in the amniotic fluid. Research is not conclusive on the relation of these factors to autism, but each of them has been identified more frequently in children with autism compared to their siblings who do not have autism and other typically developing youth. While it is unclear if any single factors during the prenatal phase affect the risk of autism, complications during pregnancy may be a risk.
There are also studies being done to test whether certain types of regressive autism have an autoimmune basis.
Maternal nutrition and inflammation during preconception and pregnancy influences fetal neurodevelopment. Intrauterine growth restriction is associated with ASD, in both term and preterm infants. Maternal inflammatory and autoimmune diseases may damage fetal tissues, aggravating a genetic problem or damaging the nervous system. Systematic reviews and meta-analyses have found that maternal prenatal infections, prenatal antibiotic exposure, and post-term pregnancies are associated with increased risk of ASD in children.
Exposure to air pollution during child pregnancy, especially heavy metals and particulates, may increase the risk of autism. Environmental factors that have been claimed without evidence to contribute to or exacerbate autism include certain foods, infectious diseases, solvents, PCBs, phthalates and phenols used in plastic products, pesticides, brominated flame retardants, alcohol, smoking, illicit drugs, vaccines, and prenatal stress. Some, such as the MMR vaccine, have been completely disproven.
Disproven vaccine hypothesis
Parents may first become aware of ASD symptoms in their child around the time of a routine vaccination. This has led to unsupported and disproven theories blaming vaccine "overload", the vaccine preservative thiomersal, or the MMR vaccine for causing autism spectrum disorder. In 1998, British physician and academic Andrew Wakefield led a fraudulent, litigation-funded study that suggested that the MMR vaccine may cause autism.
Two versions of the vaccine causation hypothesis were that autism results from brain damage caused by either the MMR vaccine itself, or by mercury used as a vaccine preservative. No convincing scientific evidence supports these claims. They are biologically implausible, and further evidence continues to refute them, including the observation that the rate of autism continues to climb despite elimination of thimerosal from most routine vaccines given to children from birth to 6 years of age.
A 2014 meta-analysis examined ten major studies on autism and vaccines involving 1.25 million children worldwide; it concluded that neither the vaccine preservative thimerosal (mercury), nor the MMR vaccine, which has never contained thimerosal, lead to the development of ASDs. Despite this, misplaced parental concern has led to lower rates of childhood immunizations, outbreaks of previously controlled childhood diseases in some countries, and the preventable deaths of several children.
Etiological hypotheses
Several hypotheses have been presented that try to explain how and why autism develops by integrating known causes (genetic and environmental effects) and findings (neurobiological and somatic). Some are more comprehensive, such as the Pathogenetic Triad, which proposes and operationalizes three core features (an autistic personality, cognitive compensation, neuropathological burden) that interact to cause autism, and the Intense World Theory, which explains autism through a hyper-active neurobiology that leads to an increased perception, attention, memory, and emotionality. There are also simpler hypotheses that explain only individual parts of the neurobiology or phenotype of autism, such as mind-blindness (a decreased ability for theory of mind), the weak central coherence theory, or the extreme male brain and empathising–systemising theory.
Evolutionary hypotheses
Research exploring the evolutionary benefits of autism and associated genes has suggested that autistic people may have played a "unique role in technological spheres and understanding of natural systems" in the course of human development. It has been suggested that autism may have arisen as "a slight trade off for other traits that are seen as highly advantageous", providing "advantages in tool making and mechanical thinking", with speculation that the condition may "reveal itself to be the result of a balanced polymorphism, like sickle cell anemia, that is advantageous in a certain mixture of genes and disadvantageous in specific combinations". In 2011, a paper in Evolutionary Psychology proposed that autistic traits, including increased spatial intelligence, concentration and memory, could have been naturally selected to enable self-sufficient foraging in a more (although not completely) solitary environment. This is called the "Solitary Forager Hypothesis". A 2016 paper examines Asperger syndrome as "an alternative prosocial adaptive strategy" that may have developed as a result of the emergence of "collaborative morality" in the context of small-scale hunter-gathering, i.e., where "a positive social reputation for making a contribution to group wellbeing and survival" becomes more important than complex social understanding.
Some research suggests that recent human evolution may be a driving force in the rise of autism in recent human populations. Studies in evolutionary medicine indicate that as cultural evolution outpaces biological evolution, disorders linked to bodily dysfunction increase in prevalence due to lack of contact with pathogens and negative environmental conditions that once widely affected ancestral populations. Because natural selection favors reproduction over health and longevity, the lack of this impetus to adapt to certain harmful circumstances creates a tendency for genes in descendant populations to over-express themselves, which may cause a wide array of maladies, ranging from mental disorders to autoimmune diseases. Conversely, noting the failure to find specific alleles that reliably cause autism or rare mutations that account for more than 5% of the heritable variation in autism established by twin and adoption studies, research in evolutionary psychiatry has concluded that it is unlikely that there is selection pressure for autism when considering that, like schizophrenics, autistic people and their siblings tend to have fewer offspring on average than non-autistic people, and instead that autism is probably better explained as a by-product of adaptive traits caused by antagonistic pleiotropy and by genes that are retained due to a fitness landscape with an asymmetric distribution.
Pathophysiology
Diagnosis
Conditions correlated or comorbid to autism
Autism is correlated or comorbid with several personality traits/disorders. Comorbidity may increase with age and may worsen the course of youth with ASDs and make intervention and treatment more difficult. Distinguishing between ASDs and other diagnoses can be challenging because the traits of ASDs often overlap with symptoms of other disorders, and the characteristics of ASDs make traditional diagnostic procedures difficult.
Correlations
Research indicates that autistic people are significantly more likely to be LGBT than the general population. There is tentative evidence that gender dysphoria occurs more frequently in autistic people. A 2021 anonymized online survey of 16- to 90-year-olds revealed that autistic males are more likely to identify as bisexual than their non-autistic peers, while autistic females are more likely to identify as homosexual than non-autistic females do.
People on the autism spectrum are significantly more likely to be non-theistic than members of the general population.
Comorbidities
The most common medical condition occurring in autistic people is seizure disorder or epilepsy, which occurs in 11–39% of autistic people. The risk varies with age, cognitive level, and type of language disorder.
Tuberous sclerosis, an autosomal dominant genetic condition in which non-malignant tumors grow in the brain and on other vital organs, is present in 1–4% of autistic people.
Intellectual disabilities are some of the most common comorbid disorders with ASDs. As diagnosis is increasingly being given to people with higher functioning autism, there is a tendency for the proportion with comorbid intellectual disability to decrease over time. In a 2019 study, it was estimated that approximately 30–40% of people diagnosed with ASD also have intellectual disability. Recent research has suggested that autistic people with intellectual disability tend to have rarer, more harmful, genetic mutations than those found in people solely diagnosed with autism. A number of genetic syndromes causing intellectual disability may also be comorbid with ASD, including fragile X, Down, Prader-Willi, Angelman, Williams syndrome, branched-chain keto acid dehydrogenase kinase deficiency, and SYNGAP1-related intellectual disability.
Learning disabilities are also highly comorbid in people with an ASD. Approximately 25–75% of people with an ASD also have some degree of a learning disability. In particular, attention deficit disorder, which is generally more prevalent than autism (ca. 8% vs. 1%), is not directly related, though it is sometimes comorbid with autism.
Various anxiety disorders tend to co-occur with ASDs, with overall comorbidity rates of 7–84%. They are common among children with ASD; there are no firm data, but studies have reported prevalences ranging from 11% to 84%. Many anxiety disorders have symptoms that are better explained by ASD itself or are hard to distinguish from ASD's symptoms.
Rates of comorbid depression in people with an ASD range from 4–58%.
The relationship between ASD and schizophrenia remains a controversial subject under continued investigation, and recent meta-analyses have examined genetic, environmental, infectious, and immune risk factors that may be shared between the two conditions. Oxidative stress, DNA damage and DNA repair have been postulated to play a role in the aetiopathology of both ASD and schizophrenia.
Deficits in ASD are often linked to behavior problems, such as difficulties following directions, being cooperative, and doing things on other people's terms. Symptoms similar to those of ADHD can be part of an ASD diagnosis.
Sensory processing disorder is also comorbid with ASD, with comorbidity rates of 42–88%.
Starting in adolescence, some people with Asperger syndrome (26% in one sample) fall under the criteria for the similar condition schizoid personality disorder, which is characterized by a lack of interest in social relationships, a tendency towards a solitary or sheltered lifestyle, secretiveness, emotional coldness, detachment and apathy. Asperger syndrome was traditionally called "schizoid disorder of childhood".
Genetic disorders – about 10–15% of autism cases have an identifiable Mendelian (single-gene) condition, chromosome abnormality, or other genetic syndromes.
Several metabolic defects, such as phenylketonuria, are associated with autistic symptoms.
Gastrointestinal problems are one of the most commonly co-occurring medical conditions in autistic people. These are linked to greater social impairment, irritability, language impairments, mood changes, and behavior and sleep problems. A 2015 review proposed that immune, gastrointestinal inflammation, malfunction of the autonomic nervous system, gut flora alterations, and food metabolites may cause brain neuroinflammation and dysfunction. A 2016 review concludes that enteric nervous system abnormalities might play a role in neurological disorders such as autism. Neural connections and the immune system are a pathway that may allow diseases originated in the intestine to spread to the brain.
Sleep problems affect about two-thirds of autistic people at some point in childhood. These most commonly include symptoms of insomnia, such as difficulty falling asleep, frequent nocturnal awakenings, and early morning awakenings. Sleep problems are associated with difficult behaviors and family stress, and are often a focus of clinical attention over and above the primary ASD diagnosis.
Dysautonomia is common in ASD, affecting heart rate and blood pressure and causing symptoms such as brain fog, blurry vision, and bowel dysfunction. It can be diagnosed through a Tilt table test.
The frequency of ASD is 10 times higher in mast cell activation syndrome patients than in the general population. This immunological condition causes cardiovascular, dermatological, gastrointestinal, neurological, and respiratory problems.
Management
There is no treatment as such for autism, and many sources advise that this is not an appropriate goal, although treatment of co-occurring conditions remains an important goal. There is no cure for autism, nor can any of the known treatments significantly reduce brain mutations caused by autism, although those who require little to no support are more likely to experience a lessening of symptoms over time. Several interventions can help children with autism, and no single treatment is best, with treatment typically tailored to the child's needs. Studies of interventions have methodological problems that prevent definitive conclusions about efficacy, but the development of evidence-based interventions has advanced.
The main goals of treatment are to lessen associated deficits and family distress, and to increase quality of life and functional independence. In general, higher IQs are correlated with greater responsiveness to treatment and improved treatment outcomes. Behavioral, psychological, education, and skill-building interventions may be used to assist autistic people to learn life skills necessary for living independently, as well as other social, communication, and language skills. Therapy also aims to reduce challenging behaviors and build upon strengths.
Intensive, sustained special education programs and behavior therapy early in life may help children acquire self-care, language, and job skills. Although evidence-based interventions for autistic children vary in their methods, many adopt a psychoeducational approach to enhancing cognitive, communication, and social skills while minimizing problem behaviors. While medications have not been found to help with core symptoms, they may be used for associated symptoms, such as irritability, inattention, or repetitive behavior patterns.
Non-pharmacological interventions
Intensive, sustained special education or remedial education programs and behavior therapy early in life may help children acquire self-care, social, and job skills. Available approaches include applied behavior analysis, developmental models, structured teaching, speech and language therapy, cognitive behavioral therapy, social skills therapy, and occupational therapy. Among these approaches, interventions either treat autistic features comprehensively, or focus treatment on a specific area of deficit. Generally, when educating those with autism, specific tactics may be used to effectively relay information to these people. Using as much social interaction as possible is key in targeting the inhibition autistic people experience concerning person-to-person contact. Additionally, research has shown that employing semantic groupings, which involves assigning words to typical conceptual categories, can be beneficial in fostering learning.
There has been increasing attention to the development of evidence-based interventions for autistic young children. Three theoretical frameworks outlined for early childhood intervention include applied behavior analysis (ABA), the developmental social-pragmatic model (DSP) and cognitive behavioral therapy (CBT). Although ABA therapy has a strong evidence base, particularly in regard to early intensive home-based therapy, ABA's effectiveness may be limited by diagnostic severity and IQ of the person affected by ASD. The Journal of Clinical Child and Adolescent Psychology has published a paper deeming two early childhood interventions "well-established": individual comprehensive ABA, and focused teacher-implemented ABA combined with DSP.
Many people have criticized ABA, calling it unhelpful and unethical. Sandoval-Norton et al. also discuss the "unintended but damaging consequences, such as prompt dependency, psychological abuse and compliance" that result in autistic people facing challenges as they transition into adulthood. Some ABA advocates have responded to such critiques that, instead of stopping ABA, there should be movement to increase protections and ethical compliance when working with autistic children.
Another evidence-based intervention that has demonstrated efficacy is a parent training model, which teaches parents how to implement various ABA and DSP techniques themselves. Various DSP programs have been developed to explicitly deliver intervention systems through at-home parent implementation.
In October 2015, the American Academy of Pediatrics (AAP) proposed new evidence-based recommendations for early interventions in ASD for children under 3. These recommendations emphasize early involvement with both developmental and behavioral methods, support by and for parents and caregivers, and a focus on both the core and associated symptoms of ASD. But a Cochrane review found no evidence that early intensive behavioral intervention (EIBI) is effective in reducing behavioral problems associated with autism in most autistic children, though it did improve IQ and language skills. The Cochrane review acknowledged that this may be due to the low quality of studies available on EIBI and therefore providers should recommend EIBI based on their clinical judgment and the family's preferences. No adverse effects of EIBI treatment were found. A meta-analysis in that same database indicates that due to the heterology in ASD, children progress to differing early intervention modalities based on ABA.
ASD treatment generally focuses on behavioral and educational interventions to target its two core symptoms: social communication deficits and restricted, repetitive behaviors. If symptoms continue after behavioral strategies have been implemented, some medications can be recommended to target specific symptoms or co-existing problems such as restricted and repetitive behaviors (RRBs), anxiety, depression, hyperactivity/inattention and sleep disturbance. Melatonin, for example, can be used for sleep problems.
Several parent-mediated behavioral therapies target social communication deficits in children with autism, but their efficacy in treating RRBs is uncertain.
Education
Educational interventions often used include applied behavior analysis (ABA), developmental models, structured teaching, speech and language therapy and social skills therapy. Among these approaches, interventions either treat autistic features comprehensively, or focalize treatment on a specific area of deficit.
The quality of research for early intensive behavioral intervention (EIBI)—a treatment procedure incorporating over 30 hours per week of the structured type of ABA that is carried out with very young children—is low; more vigorous research designs with larger sample sizes are needed. Two theoretical frameworks outlined for early childhood intervention include structured and naturalistic ABA interventions, and developmental social pragmatic models (DSP). One interventional strategy utilizes a parent training model, which teaches parents how to implement various ABA and DSP techniques, allowing for parents to disseminate interventions themselves. Various DSP programs have been developed to explicitly deliver intervention systems through at-home parent implementation. Despite the recent development of parent training models, these interventions have demonstrated effectiveness in numerous studies, being evaluated as a probable efficacious mode of treatment. Early, intensive ABA therapy has demonstrated effectiveness in enhancing communication and adaptive functioning in preschool children; it is also well-established for improving the intellectual performance of that age group.
In 2018, a Cochrane meta-analysis database concluded that some recent research is beginning to suggest that because of the heterology of ASD, there are two different ABA teaching approaches to acquiring spoken language: children with higher receptive language skills respond to 2.5 to 20 hours per week of the naturalistic approach, whereas children with lower receptive language skills require 25 hours per week of discrete trial training—the structured and intensive form of ABA. A 2023 randomized control trial study of 164 participants showed similar findings.
Similarly, a teacher-implemented intervention that utilizes a more naturalistic form of ABA combined with a developmental social pragmatic approach has been found to be beneficial in improving social-communication skills in young children, although there is less evidence in its treatment of global symptoms. Neuropsychological reports are often poorly communicated to educators, resulting in a gap between what a report recommends and what education is provided. The appropriateness of including children with varying severity of autism spectrum disorders in the general education population is a subject of current debate among educators and researchers.
Pharmacological interventions
Medications may be used to treat ASD symptoms that interfere with integrating a child into home or school when behavioral treatment fails. They may also be used for associated health problems, such as ADHD, anxiety, or if the person is hurting themself or aggressive with others, but their routine prescription for ASD's core features is not recommended. More than half of US children diagnosed with ASD are prescribed psychoactive drugs or anticonvulsants, with the most common drug classes being antidepressants, stimulants, and antipsychotics. The atypical antipsychotic drugs risperidone and aripiprazole are FDA-approved for treating associated aggressive and self-injurious behaviors. But their side effects must be weighed against their potential benefits, and autistic people may respond atypically. Side effects may include weight gain, tiredness, drooling, and paradoxical aggression. Some emerging data show positive effects of aripiprazole and risperidone on restricted and repetitive behaviors (i.e., stimming; e.g., flapping, twisting, complex whole-body movements), but due to the small sample size and different focus of these studies and the concerns about their side effects, antipsychotics are not recommended as primary treatment of RRBs. SSRI antidepressants, such as fluoxetine and fluvoxamine, have been shown to be effective in reducing repetitive and ritualistic behaviors, while the stimulant medication methylphenidate is beneficial for some children with comorbid inattentiveness or hyperactivity. There is scant reliable research about the effectiveness or safety of drug treatments for adolescents and adults with ASD. No known medication is approved for treating autism's core symptoms of social and communication impairments, although animal models indicate that postnatal administration of MDMA may be effective. MDMA has also been investigated alongside psychotherapy to treat social anxiety in autistic adults.
Alternative medicine
A multitude of alternative therapies have been researched and implemented, and many have resulted in harm to autistic people. A 2020 systematic review on adults with autism provided evidence that mindfulness-based interventions may decrease stress, anxiety, ruminating thoughts, anger, and aggression and improve mental health.
Although popularly used as an alternative treatment for autistic people, there is no good evidence to recommend a gluten- and casein-free diet as a standard treatment. A 2018 review concluded that it may be a therapeutic option for specific groups of children with autism, such as those with known food intolerances or allergies, or with food intolerance markers. The authors analyzed the prospective trials conducted to date that studied the efficacy of the gluten- and casein-free diet in children with ASD (4 in total). All of them compared gluten- and casein-free diet versus normal diet with a control group (2 double-blind randomized controlled trials, 1 double-blind crossover trial, 1 single-blind trial). In two of the studies, whose duration was 12 and 24 months, a significant improvement in ASD symptoms (efficacy rate 50%) was identified. In the other two studies, whose duration was 3 months, no significant effect was observed. The authors concluded that a longer duration of the diet may be necessary to achieve the improvement of the ASD symptoms. Other problems documented in the trials carried out include transgressions of the diet, small sample size, the heterogeneity of the participants and the possibility of a placebo effect. In the subset of people who have gluten sensitivity there is limited evidence that suggests that a gluten-free diet may improve some autistic behaviors.
The preference that autistic children have for unconventional foods can lead to reduction in bone cortical thickness with this risk being greater in those on casein-free diets, as a consequence of the low intake of calcium and vitamin D; however, suboptimal bone development in ASD has also been associated with lack of exercise and gastrointestinal disorders. In 2005, botched chelation therapy killed a five-year-old child with autism. Chelation is not recommended for autistic people since the associated risks outweigh any potential benefits. Another alternative medicine practice with no evidence is CEASE therapy, a pseudoscientific mixture of homeopathy, supplements, and "vaccine detoxing".
Results of a systematic review on interventions to address health outcomes among autistic adults found emerging evidence to support mindfulness-based interventions for improving mental health. This includes decreasing stress, anxiety, ruminating thoughts, anger, and aggression. An updated Cochrane review (2022) found evidence that music therapy likely improves social interactions, verbal communication, and nonverbal communication skills. There has been early research on hyperbaric treatments in children with autism. Studies on pet therapy have shown positive effects.
Prevention
While infection with rubella during pregnancy causes fewer than 1% of cases of autism, vaccination against rubella can prevent many of those cases.
Prognosis
There is no evidence of a cure for autism. The degree of symptoms can decrease, occasionally to the extent that people lose their diagnosis of ASD; this occurs sometimes after intensive treatment and sometimes not. It is not known how often this outcome happens, with reported rates in unselected samples ranging from 3% to 25%. Although core difficulties tend to persist, symptoms often become less severe with age. Acquiring language before age six, having an IQ above 50, and having a marketable skill all predict better outcomes; independent living is unlikely in autistic people with higher support needs.
Among others, academic Temple Grandin has advised against striving to cure autism, saying that if a cure were found, she would choose to stay the way she is. She wrote, "The skills that people with autism bring to the table should be nurtured for their benefit and [for the benefit of] society", adding, "If you totally get rid of autism, you'd have nobody to fix your computer in the future".
The prognosis of autism describes the developmental course, gradual autism development, regressive autism development, differential outcomes, academic performance and employment.
Epidemiology
The World Health Organization estimates about 1 in 100 children had autism during the period from 2012 to 2021 as that was the average estimate in studies published during that period with a trend of increasing prevalence over time. However, the study's 1% figure may reflect an underestimate of prevalence in low- and middle-income countries. The number of people diagnosed has increased considerably since the 1990s, which may be partly due to increased recognition of the condition.
While rates of ASD are consistent across cultures, they vary greatly by gender, with boys diagnosed far more frequently than girls: 1 in 70 boys, but only 1 in 315 girls at eight years of age. Girls, however, are more likely to have associated cognitive impairment, suggesting that less severe forms of ASD are likely being missed in girls and women. Prevalence differences may be a result of gender differences in expression of clinical symptoms, with women and girls with autism showing less atypical behaviors and, therefore, less likely to receive an ASD diagnosis.
Using DSM-5 criteria, 92% of the children diagnosed per DSM-IV with one of the disorders which is considered part of ASD will still meet the diagnostic criteria of ASD. However, if both ASD and the social (pragmatic) communication disorder categories of DSM-5 are combined, the prevalence of autism is mostly unchanged from the prevalence per the DSM-IV criteria. The best estimate for prevalence of ASD is 0.7% or 1 child in 143 children. Relatively mild forms of autism, such as Asperger's as well as other developmental disorders, are included in the DSM-5 diagnostic criteria. ASD rates were constant between 2014 and 2016 but twice the rate compared to the time period between 2011 and 2014 (1.25 vs 2.47%). A Canadian meta-analysis from 2019 confirmed these effects as the profiles of autistic people became less and less different from the profiles of the general population. In the US, the rates for diagnosed ASD have been steadily increasing since 2000 when records began being kept. While it remains unclear whether this trend represents a true rise in incidence, it likely reflects changes in ASD diagnostic criteria, improved detection, and increased public awareness of autism. In 2012, the NHS estimated that the overall prevalence of autism among adults aged 18 years and over in the UK was 1.1%. A 2016 survey in the United States reported a rate of 25 per 1,000 children for ASD. Rates of autism are poorly understood in many low- and middle-income countries, which affects the accuracy of global ASD prevalence estimates, but it is thought that most autistic people live in low- and middle-income countries.
In 2020, the Centers for Disease Control's Autism and Developmental Disabilities Monitoring (ADDM) Network reported that approximately 1 in 54 children in the United States (1 in 34 boys, and 1 in 144 girls) is diagnosed with an autism spectrum disorder (ASD), based on data collected in 2016. This estimate is a 10% increase from the 1 in 59 rate in 2014, a 105% increase from the 1 in 110 rate in 2006, and a 176% increase from the 1 in 150 rate in 2000. Diagnostic criteria for ASD have changed significantly since the 1980s; for example, U.S. special-education autism classification was introduced in 1994.
In the UK, from 1998 to 2018, the autism diagnoses increased by 787%. This increase is largely attributable to changes in diagnostic practices, referral patterns, availability of services, age at diagnosis, and public awareness, particularly among women, though unidentified environmental risk factors cannot be ruled out. The available evidence does not rule out the possibility that autism's true prevalence has increased; a real increase would suggest directing more attention and funding toward psychosocial factors and changing environmental factors instead of continuing to focus on genetics. It has been established that vaccination is not a risk factor for autism and is not a cause of any increase in autism prevalence rates, if any change in the rate of autism exists at all.
Males have higher likelihood of being diagnosed with ASD than females. The sex ratio averages 4.3:1 and is greatly modified by cognitive impairment: it may be close to 2:1 with intellectual disability and more than 5.5:1 without. Several theories about the higher prevalence in males have been investigated, but the cause of the difference is unconfirmed; one theory is that females are underdiagnosed.
The risk of developing autism is greater with older fathers than with older mothers; two potential explanations are the known increase in mutation burden in older sperm, and the hypothesis that men marry later if they carry genetic liability and show some signs of autism. Most professionals believe that race, ethnicity, and socioeconomic background do not affect the occurrence of autism.
History
Society and culture
An autistic culture has emerged, accompanied by the autistic rights and neurodiversity movements, that argues autism should be accepted as a difference to be accommodated instead of cured, although a minority of autistic people might still accept a cure. Worldwide, events related to autism include World Autism Awareness Day, Autism Sunday, Autistic Pride Day, Autreat, and others.
Social-science scholars study those with autism in hopes to learn more about "autism as a culture, transcultural comparisons ... and research on social movements." Many autistic people have been successful in their fields.
Special interests are commonly found in autistic people, sometimes leading to hobbies, vast collections, and activism. Environmental activist Greta Thunberg has spoken favorably about her autism diagnosis, saying that autism can be a source of life purpose, as well as forming the basis of careers, hobbies, and friendships.
Neurodiversity movement
Some autistic people, as well as a growing number of researchers, have advocated a shift in attitudes toward the view that autism spectrum disorder is a difference, rather than a disease that must be treated or cured. Critics have bemoaned the entrenchment of some of these groups' opinions.
The neurodiversity movement and the autism rights movement are social movements within the context of disability rights, emphasizing the concept of neurodiversity, which describes the autism spectrum as a result of natural variations in the human brain rather than a disorder to be cured. The autism rights movement advocates including greater acceptance of autistic behaviors, therapies that focus on coping skills rather than imitating the behaviors of those without autism, and the recognition of the autistic community as a minority group.
Autism rights or neurodiversity advocates believe that the autism spectrum is genetic and should be accepted as a natural variation in the human genome. These movements are not without detractors; a common argument against neurodiversity activists is that most of them have relatively low support needs, or are self-diagnosed, and do not represent the views of autistic people with higher support needs. Jacquiline den Houting explores this critique, determining that the voices of low-support needs autistics are "some of the most influential within the neurodiversity movement, although admittedly these voices are a minority within the advocacy community"; she suggests this is in part a shortcoming of the wider neurotypical community, referencing nonspeaking self-advocate Amy Sequenzia's writing. Pier Jaarsma and Stellan Welin make the argument that only high-functioning autistic people should be included under the neurodiversity banner, as low-functioning autists' condition may rightfully be viewed as a disability. The concept of neurodiversity is contentious in autism advocacy and research groups and has led to infighting.
Events
Since 2011, the Autistic Self Advocacy Network has celebrated April as Autism Acceptance Month. In 2021, the Autism Society of America urged organizations to retitle Autism Awareness Day as Autism Acceptance Day, to focus on "more fully integrating those 1 in 54 Americans living with autism into our social fabric".
Symbols and flags
Puzzle piece
In 1963, the British National Autistic Society chose a puzzle piece as its logo, due to its view of autistic people as suffering from a "puzzling" condition. The logo, designed by board member Gerald Gasson, consisted of a green and black puzzle piece with four knobs, with a crying child at its center. Other organizations and advocates adopted the puzzle piece as a symbol of autism, including American organization Autism Speaks, which uses a puzzle piece with one knob, two holes, one edge.
In 1999, the Autism Society designed the puzzle ribbon (an awareness ribbon patterned with red, yellow, cyan, and blue puzzle pieces) as a symbol of autism awareness.
The puzzle symbol is controversial among autism advocates and rejected by many. It has been criticized as outdated, now that autism is better understood, as well as implying that autistic people are mysterious or incomplete, and for its association with Autism Speaks. The autism rights movement and neurodiversity advocates have criticized Autism Speaks for its view of autism as a disease to be cured.
Rainbow infinity
In 2004, neurodiversity advocates Amy and Gwen Nelson designed the "rainbow infinity symbol", originally as the logo for their advocacy group Aspies For Freedom. Many adopted the infinity symbol as a symbol for the autism spectrum. The prismatic colors are often associated with the neurodiversity movement in general.
In 2018, Julian Morgan wrote the article "Light It Up Gold", a response to the "Light It Up Blue" awareness campaign Autism Speaks launched in 2007. Morgan pushed to use gold to symbolize autism, due its chemical symbol Au, from the Latin .
Flags
An autistic pride flag was created in 2005 by Aspies For Freedom for the first Autistic Pride Day, featuring a rainbow infinity symbol on a white background.
As the rainbow infinity on a white background has become increasingly viewed as representative of neurodiversity in general, several designs have been proposed for an autistic-specific flag. In 2023, the People's History Museum featured a 2015 autistic pride design by Joseph Redford, featuring a rainbow infinity symbol, a green background for being true to one's nature, and a purple background for neurodiversity.
Caregivers
Families who care for an autistic child face added stress from a number of different causes. Parents may struggle to understand the diagnosis and to find appropriate care options. They often take a negative view of the diagnosis, and may struggle emotionally. More than half of parents over age 50 are still living with their child, as about 85% of autistic people have difficulties living independently. Some studies also find decreased earnings among parents who care for autistic children. Siblings of children with ASD report greater admiration and less conflict with the affected sibling than siblings of unaffected children, like siblings of children with Down syndrome. But they reported lower levels of closeness and intimacy than siblings of children with Down syndrome; siblings of autistic people have a greater risk of negative well-being and poorer sibling relationships as adults.
| Biology and health sciences | Mental disorder | null |
945670 | https://en.wikipedia.org/wiki/Cirque | Cirque | A (; from the Latin word ) is an amphitheatre-like valley formed by glacial erosion. Alternative names for this landform are corrie (from , meaning a pot or cauldron) and ; ). A cirque may also be a similarly shaped landform arising from fluvial erosion.
The concave shape of a glacial cirque is open on the downhill side, while the cupped section is generally steep. Cliff-like slopes, down which ice and glaciated debris combine and converge, form the three or more higher sides. The floor of the cirque ends up bowl-shaped, as it is the complex convergence zone of combining ice flows from multiple directions and their accompanying rock burdens. Hence, it experiences somewhat greater erosion forces and is most often overdeepened below the level of the cirque's low-side outlet (stage) and its down-slope (backstage) valley. If the cirque is subject to seasonal melting, the floor of the cirque most often forms a tarn (small lake) behind a dam, which marks the downstream limit of the glacial overdeepening. The dam itself can be composed of moraine, glacial till, or a lip of the underlying bedrock.
The fluvial cirque or , found in karst landscapes, is formed by intermittent river flow cutting through layers of limestone and chalk leaving sheer cliffs. A common feature for all fluvial-erosion cirques is a terrain which includes erosion resistant upper structures overlying materials which are more easily eroded.
Formation
Glacial-erosion cirque formation
Glacial cirques are found amongst mountain ranges throughout the world; 'classic' cirques are typically about one kilometer long and one kilometer wide. Situated high on a mountainside near the firn line, they are typically partially surrounded on three sides by steep cliffs. The highest cliff is often called a headwall. The fourth side forms the lip, threshold or sill, the side at which the glacier flowed away from the cirque. Many glacial cirques contain tarns dammed by either till (debris) or a bedrock threshold. When enough snow accumulates, it can flow out the opening of the bowl and form valley glaciers which may be several kilometers long.
Cirques form in conditions which are favorable; in the Northern Hemisphere the conditions include the north-east slope, where they are protected from the majority of the Sun's energy and from the prevailing winds. These areas are sheltered from heat, encouraging the accumulation of snow; if the accumulation of snow increases, the snow turns into glacial ice. The process of nivation follows, whereby a hollow in a slope may be enlarged by ice segregation weathering and glacial erosion. Ice segregation erodes the vertical rock face and causes it to disintegrate, which may result in an avalanche bringing down more snow and rock to add to the growing glacier. Eventually, this hollow may become large enough that glacial erosion intensifies. The enlarging of this open ended concavity creates a larger leeward deposition zone, furthering the process of glaciation. Debris (or till) in the ice also may abrade the bed surface; should ice move down a slope it would have a 'sandpaper effect' on the bedrock beneath, on which it scrapes.
Eventually, the hollow may become a large bowl shape in the side of the mountain, with the headwall being weathered by ice segregation, and as well as being eroded by plucking. The basin will become deeper as it continues to be eroded by ice segregation and abrasion. Should ice segregation, plucking and abrasion continue, the dimensions of the cirque will increase, but the proportion of the landform would remain roughly the same. A bergschrund forms when the movement of the glacier separates the moving ice from the stationary ice, forming a crevasse. The method of erosion of the headwall lying between the surface of the glacier and the cirque's floor has been attributed to freeze-thaw mechanisms. The temperature within the bergschrund changes very little, however, studies have shown that ice segregation (frost shattering) may happen with only small changes in temperature. Water that flows into the bergschrund can be cooled to freezing temperatures by the surrounding ice, allowing freeze-thaw free mechanisms to occur.
If two adjacent cirques erode toward one another, an arête, or steep sided ridge, forms. When three or more cirques erode toward one another, a pyramidal peak is created. In some cases, this peak will be made accessible by one or more arêtes. The Matterhorn in the European Alps is an example of such a peak.
Where cirques form one behind the other, a cirque stairway results, as at the Zastler Loch in the Black Forest.
As glaciers can only originate above the snowline, studying the location of present-day cirques provides information on past glaciation patterns and on climate change.
Fluvial-erosion cirque formation
Although a less common usage, the term cirque is also used for amphitheatre-shaped, fluvial-erosion features. For example, an approximately anticlinal erosion cirque is at on the southern boundary of the Negev highlands. This erosional cirque or was formed by intermittent river flow in the Makhtesh Ramon cutting through layers of limestone and chalk, resulting in cirque walls with a sheer drop. The Cirque du Bout du Monde is another such feature, created in karst terraine in the Burgundy region of the department of in France.
Yet another type of fluvial erosion-formed cirque is found on Réunion island, which includes the tallest volcanic structure in the Indian Ocean. The island consists of an active shield-volcano () and an extinct, deeply eroded volcano (Piton des Neiges). Three cirques have eroded there in a sequence of agglomerated, fragmented rock and volcanic breccia associated with pillow lavas overlain by more coherent, solid lavas.
A common feature for all fluvial-erosion cirques is a terrain which includes erosion resistant upper structures overlying materials which are more easily eroded.
Notable cirques
Australia
Blue Lake Cirque, New South Wales, Australia
Asia
Chandra Taal, Himachal Pradesh, India
Cirque Valley, Hindu Kush, Pakistan
Karasawa Cirque, Kamikōchi, Mount Hotakadake, Hida Mountains, Japan
Makhtesh Ramon, Negev desert, Israel
Senjōjiki Cirque, Mount Hōken, Kiso Mountains, Japan
Western Cwm, Khumbu Himal, Nepal
Yamasaki Cirque, Mount Tateyama, Japan
Europe (glacial)
Cadair Idris, Wales
Circo de Gredos, Sierra de Gredos, Spain
Cirque de Gavarnie, Pyrenees, France
Cirque d'Estaubé, Pyrenees, France
Maritsa cirque, Rila Mountain, Bulgaria
Malyovitsa cirque, Rila Mountain, Bulgaria
Seven Rila Lakes cirques, Rila Mountain, Bulgaria
Banderishki cirque, Pirin Mountain, Bulgaria
, Grampian Mountains, Scottish Highlands
Śnieżne Kotły, Karkonosze, Poland
Coumshingaun Lake, County Waterford, Ireland
Europe (fluvial)
Cirque de Navacelles, Grands Causses, France
Cirque du Bout du Monde, Grands Causses, France
Cirque du Bout du Monde, Burgundy, France
North America
Cirque of the Towers, Wyoming, United States
Iceberg Cirque, Montana, US
Summit Lake cirque, and others on Mount Blue Sky, Colorado, US
Great Basin and others on Mount Katahdin, Maine, US
Great Gulf, New Hampshire, US
Tuckerman Ravine, New Hampshire, US
| Physical sciences | Glacial landforms | Earth science |
945696 | https://en.wikipedia.org/wiki/Acanthodii | Acanthodii | Acanthodii or acanthodians is an extinct class of gnathostomes (jawed fishes). They are currently considered to represent a paraphyletic grade of various fish lineages basal to extant Chondrichthyes, which includes living sharks, rays, and chimaeras. Acanthodians possess a mosaic of features shared with both osteichthyans (bony fish) and chondrichthyans (cartilaginous fish). In general body shape, they were similar to modern sharks, but their epidermis was covered with tiny rhomboid platelets like the scales of holosteians (gars, bowfins).
The popular name "spiny sharks" is because they were superficially shark-shaped, with a streamlined body, paired fins, a strongly upturned tail, and stout, largely immovable bony spines supporting all the fins except the tail—hence, "spiny sharks". However, acanthodians are not true sharks; their close relation to modern cartilaginous fish can lead them to be considered "stem-sharks". Acanthodians had a cartilaginous skeleton, but their fins had a wide, bony base and were reinforced on their anterior margin with a dentine spine. As a result, fossilized spines and scales are often all that remains of these fishes in ancient sedimentary rocks. The earliest acanthodians were marine, but during the Devonian, freshwater species became predominant.
Acanthodians have been divided into four orders: Acanthodiformes, Climatiiformes, Diplacanthiformes, and Ischnacanthiformes. "Climatiiformes" is a paraphyletic assemblage of early acanthodians such as climatiids, gyracanthids, and diplacanthids; they had robust bony shoulder girdles and many small sharp spines ("intermediate" or "prepelvic" spines) between the pectoral and pelvic fins. The climatiiform subgroup Diplacanthida has subsequently been elevated to its own order, Diplacanthiformes. Ischnacanthiforms were predators with tooth plates fused to their jaws. Acanthodiforms were filter feeders with a single dorsal fin, toothless jaws, and long gill rakers. They were the last and most specialized off the traditional acanthodians, as they survived up until the Permian period.
Characteristics
The scales of Acanthodii have distinctive ornamentation peculiar to each order. Because of this, the scales are often used in determining relative age of sedimentary rock. The scales are tiny, with a bulbous base, a neck, and a flat or slightly curved diamond-shaped crown.
Despite being called "spiny sharks", acanthodians predate sharks. Scales that have been tentatively identified as belonging to acanthodians, or "shark-like fishes" have been found in various Ordovician strata, though, they are ambiguous, and may actually belong to jawless fishes such as thelodonts. The earliest unequivocal acanthodian fossils date from the beginning of the Silurian Period, some 50 million years before the first sharks appeared. Later, the acanthodians colonized fresh waters, and thrived in the rivers and lakes during the Devonian and in the coal swamps of Carboniferous. By this time bony fishes were already showing their potential to dominate the waters of the world, and their competition proved too much for the spiny sharks, which died out in Permian times (approximately 250 million years ago).
Many palaeontologists originally considered the acanthodians close to the ancestors of the bony fishes. Although their interior skeletons were made of cartilage, a bonelike material had developed in the skins of these fishes, in the form of closely fitting scales (see above). Some scales were greatly enlarged and formed a bony covering on top of the head and over the lower shoulder girdle. Others developed a bony flap over the gill openings analogous to the operculum in later bony fishes. However, most of these characteristics are considered homologous characteristics derived from common placoderm ancestors, and present also in basal cartilaginous fish. Overall, the acanthodians' jaws are presumed to have evolved from the first gill arch of some ancestral jawless fishes that had a gill skeleton made of pieces of jointed cartilage.
Taxonomy and phylogeny
In a study of early jawed vertebrate relationships, Davis et al. (2012) found acanthodians to be split among the two major clades Osteichthyes (bony fish) and Chondrichthyes (cartilaginous fish). The well-known acanthodian Acanthodes was placed within Osteichthyes, despite the presence of many chondrichthyan characteristics in its braincase. However, a newly described Silurian placoderm, Entelognathus, which has jaw anatomy shared with bony fish and tetrapods, has led to revisions of this phylogeny: acanthodians were then considered to be a paraphyletic assemblage leading to cartilaginous fish, while bony fish evolved from placoderm ancestors.
Burrow et al. 2016 provides vindication by finding chondrichthyans to be nested among Acanthodii, most closely related to Doliodus and Tamiobatis. A 2017 study of Doliodus morphology points out that it appears to display a mosaic of shark and acanthodian features, making it a transitional fossil and further reinforcing this idea.
Phylogeny after
Evolutionary history
The oldest remains attributed acanthodian-grade chondrichthyans are Fanjingshania and Qianodus from the Early Silurian of China, dating to around 439 million years ago. Compared to other contemporary groups of fish, acanthodians were relatively morphologically and ecologically conservative. Acanthodians rose in diversity during the Late Silurian, reaching their apex of diversity during the Lochkovian stage of the Early Devonian, declining during the Pragian but rising again during the following Emsian, which was followed by a decline in diversity during middle-Late Devonian. The diversity of the group was consistently low but stable during the Carboniferous, slightly decreasing going into the Permian. The youngest records of the group are isolated scales and fin spines from Middle-Late Permian strata in the Paraná Basin of Brazil.
| Biology and health sciences | Fishes | null |
945957 | https://en.wikipedia.org/wiki/New%20Foundations | New Foundations | In mathematical logic, New Foundations (NF) is a non-well-founded, finitely axiomatizable set theory conceived by Willard Van Orman Quine as a simplification of the theory of types of Principia Mathematica.
Definition
The well-formed formulas of NF are the standard formulas of propositional calculus with two primitive predicates equality () and membership (). NF can be presented with only two axiom schemata:
Extensionality: Two objects with the same elements are the same object; formally, given any set A and any set B, if for every set X, X is a member of A if and only if X is a member of B, then A is equal to B.
A restricted axiom schema of comprehension: exists for each stratified formula .
A formula is said to be stratified if there exists a function f from pieces of 's syntax to the natural numbers, such that for any atomic subformula of we have f(y) = f(x) + 1, while for any atomic subformula of , we have f(x) = f(y).
Finite axiomatization
NF can be finitely axiomatized. One advantage of such a finite axiomatization is that it eliminates the notion of stratification. The axioms in a finite axiomatization correspond to natural basic constructions, whereas stratified comprehension is powerful but not necessarily intuitive. In his introductory book, Holmes opted to take the finite axiomatization as basic, and prove stratified comprehension as a theorem. The precise set of axioms can vary, but includes most of the following, with the others provable as theorems:
Extensionality: If and are sets, and for each object , is an element of if and only if is an element of , then . This can also be viewed as defining the equality symbol.
Singleton: For every object , the set exists, and is called the singleton of .
Cartesian Product: For any sets , , the set , called the Cartesian product of and , exists. This can be restricted to the existence of one of the cross products or .
Converse: For each relation , the set exists; observe that exactly if .
Singleton Image: For any relation , the set , called the singleton image of , exists.
Domain: If is a relation, the set , called the domain of , exists. This can be defined using the operation of type lowering.
Inclusion: The set exists. Equivalently, we may consider the set .
Complement: For each set , the set , called the complement of , exists.
(Boolean) Union: If and are sets, the set , called the (Boolean) union of and , exists.
Universal Set: exists. It is straightforward that for any set , .
Ordered Pair: For each , , the ordered pair of and , , exists; exactly if and . This and larger tuples can be a definition rather than an axiom if an ordered pair construction is used.
Projections: The sets and exist (these are the relations which an ordered pair has to its first and second terms, which are technically referred to as its projections).
Diagonal: The set exists, called the equality relation.
Set Union: If is a set all of whose elements are sets, the set , called the (set) union of , exists.
Relative Product: If , are relations, the set , called the relative product of and , exists.
Anti-intersection: exists. This is equivalent to complement and union together, with and .
Cardinal one: The set of all singletons, , exists.
Tuple Insertions: For a relation , the sets and exist.
Type lowering: For any set , the set exists.
Typed Set Theory
New Foundations is closely related to Russellian unramified typed set theory (TST), a streamlined version of the theory of types of Principia Mathematica with a linear hierarchy of types. In this many-sorted theory, each variable and set is assigned a type. It is customary to write the type indices as superscripts: denotes a variable of type n. Type 0 consists of individuals otherwise undescribed. For each (meta-) natural number n, type n+1 objects are sets of type n objects; objects connected by identity have equal types and sets of type n have members of type n-1. The axioms of TST are extensionality, on sets of the same (positive) type, and comprehension, namely that if is a formula, then the set exists. In other words, given any formula , the formula is an axiom where represents the set and is not free in . This type theory is much less complicated than the one first set out in the Principia Mathematica, which included types for relations whose arguments were not necessarily all of the same types.
There is a correspondence between New Foundations and TST in terms of adding or erasing type annotations. In NF's comprehension schema, a formula is stratified exactly when the formula can be assigned types according to the rules of TST. This can be extended to map every NF formula to a set of corresponding TST formulas with various type index annotations. The mapping is one-to-many because TST has many similar formulas. For example, raising every type index in a TST formula by 1 results in a new, valid TST formula.
Tangled Type Theory
Tangled Type Theory (TTT) is an extension of TST where each variable is typed by an ordinal rather than a natural number. The well-formed atomic formulas are and where . The axioms of TTT are those of TST where each variable of type is mapped to a variable where is an increasing function.
TTT is considered a "weird" theory because each type is related to each lower type in the same way. For example, type 2 sets have both type 1 members and type 0 members, and extensionality axioms assert that a type 2 set is determined uniquely by either its type 1 members or its type 0 members. Whereas TST has natural models where each type is the power set of type , in TTT each type is being interpreted as the power set of each lower type simultaneously. Regardless, a model of NF can be easily converted to a model of TTT, because in NF all the types are already one and the same. Conversely, with a more complicated argument, it can also be shown that the consistency of TTT implies the consistency of NF.
NFU and other variants
NF with urelements (NFU) is an important variant of NF due to Jensen and clarified by Holmes. Urelements are objects that are not sets and do not contain any elements, but can be contained in sets. One of the simplest forms of axiomatization of NFU regards urelements as multiple, unequal empty sets, thus weakening the extensionality axiom of NF to:
Weak extensionality: Two non-empty objects with the same elements are the same object; formally,
In this axiomatization, the comprehension schema is unchanged, although the set will not be unique if it is empty (i.e. if is unsatisfiable).
However, for ease of use, it is more convenient to have a unique, "canonical" empty set. This can be done by introducing a sethood predicate to distinguish sets from atoms. The axioms are then:
Sets: Only sets have members, i.e.
Extensionality: Two sets with the same elements are the same set, i.e.
Comprehension: The set exists for each stratified formula , i.e.
NF3 is the fragment of NF with full extensionality (no urelements) and those instances of comprehension which can be stratified using at most three types. NF4 is the same theory as NF.
Mathematical Logic (ML) is an extension of NF that includes proper classes as well as sets. ML was proposed by Quine and revised by Hao Wang, who proved that NF and the revised ML are equiconsistent.
Constructions
This section discusses some problematic constructions in NF. For a further development of mathematics in NFU, with a comparison to the development of the same in ZFC, see implementation of mathematics in set theory.
Ordered pairs
Relations and functions are defined in TST (and in NF and NFU) as sets of ordered pairs in the usual way. For purposes of stratification, it is desirable that a relation or function is merely one type higher than the type of the members of its field. This requires defining the ordered pair so that its type is the same as that of its arguments (resulting in a type-level ordered pair). The usual definition of the ordered pair, namely , results in a type two higher than the type of its arguments a and b. Hence for purposes of determining stratification, a function is three types higher than the members of its field. NF and related theories usually employ Quine's set-theoretic definition of the ordered pair, which yields a type-level ordered pair. However, Quine's definition relies on set operations on each of the elements a and b, and therefore does not directly work in NFU.
As an alternative approach, Holmes takes the ordered pair (a, b) as a primitive notion, as well as its left and right projections and , i.e., functions such that and (in Holmes' axiomatization of NFU, the comprehension schema that asserts the existence of for any stratified formula is considered a theorem and only proved later, so expressions like are not considered proper definitions). Fortunately, whether the ordered pair is type-level by definition or by assumption (i.e., taken as primitive) usually does not matter.
Natural numbers and the axiom of infinity
The usual form of the axiom of infinity is based on the von Neumann construction of the natural numbers, which is not suitable for NF, since the description of the successor operation (and many other aspects of von Neumann numerals) is necessarily unstratified. The usual form of natural numbers used in NF follows Frege's definition, i.e., the natural number n is represented by the set of all sets with n elements. Under this definition, 0 is easily defined as , and the successor operation can be defined in a stratified way: Under this definition, one can write down a statement analogous to the usual form of the axiom of infinity. However, that statement would be trivially true, since the universal set would be an inductive set.
Since inductive sets always exist, the set of natural numbers can be defined as the intersection of all inductive sets. This definition enables mathematical induction for stratified statements , because the set can be constructed, and when satisfies the conditions for mathematical induction, this set is an inductive set.
Finite sets can then be defined as sets that belong to a natural number. However, it is not trivial to prove that is not a "finite set", i.e., that the size of the universe is not a natural number. Suppose that . Then (it can be shown inductively that a finite set is not equinumerous with any of its proper subsets), , and each subsequent natural number would be too, causing arithmetic to break down. To prevent this, one can introduce the axiom of infinity for NF:
It may intuitively seem that one should be able to prove Infinity in NF(U) by constructing any "externally" infinite sequence of sets, such as . However, such a sequence could only be constructed through unstratified constructions (evidenced by the fact that TST itself has finite models), so such a proof could not be carried out in NF(U). In fact, Infinity is logically independent of NFU: There exists models of NFU where is a non-standard natural number. In such models, mathematical induction can prove statements about , making it impossible to "distinguish" from standard natural numbers.
However, there are some cases where Infinity can be proven (in which cases it may be referred to as the theorem of infinity):
In NF (without urelements), Specker has shown that the axiom of choice is false. Since it can be proved through induction that every finite set has a choice function (a stratified condition), it follows that is infinite.
In NFU with axioms asserting the existence of a type-level ordered pair, is equinumerous with its proper subset , which implies Infinity. Conversely, NFU + Infinity + Choice proves the existence of a type-level ordered pair. NFU + Infinity interprets NFU + "there is a type-level ordered pair" (they are not quite the same theory, but the differences are inessential).
Stronger axioms of infinity exist, such as that the set of natural numbers is a strongly Cantorian set, or NFUM = NFU + Infinity + Large Ordinals + Small Ordinals which is equivalent to Morse–Kelley set theory plus a predicate on proper classes which is a κ-complete nonprincipal ultrafilter on the proper class ordinal κ.
Large sets
NF (and NFU + Infinity + Choice, described below and known consistent) allow the construction of two kinds of sets that ZFC and its proper extensions disallow because they are "too large" (some set theories admit these entities under the heading of proper classes):
The universal set V. Because is a stratified formula, the universal set V = {x | x=x} exists by Comprehension. An immediate consequence is that all sets have complements, and the entire set-theoretic universe under NF has a Boolean structure.
Cardinal and ordinal numbers. In NF (and TST), the set of all sets having n elements (the circularity here is only apparent) exists. Hence Frege's definition of the cardinal numbers works in NF and NFU: a cardinal number is an equivalence class of sets under the relation of equinumerosity: the sets A and B are equinumerous if there exists a bijection between them, in which case we write . Likewise, an ordinal number is an equivalence class of well-ordered sets.
Cartesian closure
The category whose objects are the sets of NF and whose arrows are the functions between those sets is not Cartesian closed; Since NF lacks Cartesian closure, not every function curries as one might intuitively expect, and NF is not a topos.
Resolution of set-theoretic paradoxes
NF may seem to run afoul of problems similar to those in naive set theory, but this is not the case. For example, the existence of the impossible Russell class is not an axiom of NF, because cannot be stratified. NF steers clear of the three well-known paradoxes of set theory in drastically different ways than how those paradoxes are resolved in well-founded set theories such as ZFC. Many useful concepts that are unique to NF and its variants can be developed from the resolution of those paradoxes.
Russell's paradox
The resolution of Russell's paradox is trivial: is not a stratified formula, so the existence of is not asserted by any instance of Comprehension. Quine said that he constructed NF with this paradox uppermost in mind.
Cantor's paradox and Cantorian sets
Cantor's paradox boils down to the question of whether there exists a largest cardinal number, or equivalently, whether there exists a set with the largest cardinality. In NF, the universal set is obviously a set with the largest cardinality. However, Cantor's theorem says (given ZFC) that the power set of any set is larger than (there can be no injection (one-to-one map) from into ), which seems to imply a contradiction when .
Of course there is an injection from into since is the universal set, so it must be that Cantor's theorem (in its original form) does not hold in NF. Indeed, the proof of Cantor's theorem uses the diagonalization argument by considering the set . In NF, and should be assigned the same type, so the definition of is not stratified. Indeed, if is the trivial injection , then is the same (ill-defined) set in Russell's paradox.
This failure is not surprising since makes no sense in TST: the type of is one higher than the type of . In NF, is a syntactical sentence due to the conflation of all the types, but any general proof involving Comprehension is unlikely to work.
The usual way to correct such a type problem is to replace with , the set of one-element subsets of . Indeed, the correctly typed version of Cantor's theorem is a theorem in TST (thanks to the diagonalization argument), and thus also a theorem in NF. In particular, : there are fewer one-element sets than sets (and so fewer one-element sets than general objects, if we are in NFU). The "obvious" bijection from the universe to the one-element sets is not a set; it is not a set because its definition is unstratified. Note that in all models of NFU + Choice it is the case that ; Choice allows one not only to prove that there are urelements but that there are many cardinals between and .
However, unlike in TST, is a syntactical sentence in NF(U), and as shown above one can talk about its truth value for specific values of (e.g. when it is false). A set which satisfies the intuitively appealing is said to be Cantorian: a Cantorian set satisfies the usual form of Cantor's theorem. A set which satisfies the further condition that , the restriction of the singleton map to A, is a set is not only Cantorian set but strongly Cantorian.
Burali-Forti paradox and the T operation
The Burali-Forti paradox of the largest ordinal number is resolved in the opposite way: In NF, having access to the set of ordinals does not allow one to construct a "largest ordinal number". One can construct the ordinal that corresponds to the natural well-ordering of all ordinals, but that does not mean that is larger than all those ordinals.
To formalize the Burali-Forti paradox in NF, it is necessary to first formalize the concept of ordinal numbers. In NF, ordinals are defined (in the same way as in naive set theory) as equivalence classes of well-orderings under isomorphism. This is a stratified definition, so the set of ordinals can be defined with no problem. Transfinite induction works on stratified statements, which allows one to prove that the natural ordering of ordinals ( iff there exists well-orderings such that is a continuation of ) is a well-ordering of . By definition of ordinals, this well-ordering also belongs to an ordinal . In naive set theory, one would go on to prove by transfinite induction that each ordinal is the order type of the natural order on the ordinals less than , which would imply an contradiction since by definition is the order type of all ordinals, not any proper initial segment of them.
However, the statement " is the order type of the natural order on the ordinals less than " is not stratified, so the transfinite induction argument does not work in NF. In fact, "the order type of the natural order on the ordinals less than " is at least two types higher than : The order relation is one type higher than assuming that is a type-level ordered pair, and the order type (equivalence class) is one type higher than . If is the usual Kuratowski ordered pair (two types higher than and ), then would be four types higher than .
To correct such a type problem, one needs the T operation, , that "raises the type" of an ordinal , just like how "raises the type" of the set . The T operation is defined as follows: If , then is the order type of the order . Now the lemma on order types may be restated in a stratified manner:
The order type of the natural order on the ordinals is or , depending on which ordered pair is used.
Both versions of this statement can be proven by transfinite induction; we assume the type level pair hereinafter. This means that is always less than , the order type of all ordinals. In particular, .
Another (stratified) statement that can be proven by transfinite induction is that T is a strictly monotone (order-preserving) operation on the ordinals, i.e., iff . Hence the T operation is not a function: The collection of ordinals cannot have a least member, and thus cannot be a set. More concretely, the monotonicity of T implies , a "descending sequence" in the ordinals which also cannot be a set.
One might assert that this result shows that no model of NF(U) is "standard", since the ordinals in any model of NFU are externally not well-ordered. This is a philosophical question, not a question of what can be proved within the formal theory. Note that even within NFU it can be proven that any set model of NFU has non-well-ordered "ordinals"; NFU does not conclude that the universe is a model of NFU, despite being a set, because the membership relation is not a set relation.
Consistency
Some mathematicians have questioned the consistency of NF, partly because it is not clear why it avoids the known paradoxes. A key issue was that Specker proved NF combined with the Axiom of Choice is inconsistent. The proof is complex and involves T-operations. However, since 2010, Holmes has claimed to have shown that NF is consistent relative to the consistency of standard set theory (ZFC). In 2024, Sky Wilshaw confirmed Holmes' proof using the Lean proof assistant.
Although NFU resolves the paradoxes similarly to NF, it has a much simpler consistency proof. The proof can be formalized within Peano Arithmetic (PA), a theory weaker than ZF that most mathematicians accept without question. This does not conflict with Gödel's second incompleteness theorem because NFU does not include the Axiom of Infinity and therefore PA cannot be modeled in NFU, avoiding a contradiction. PA also proves that NFU with Infinity and NFU with both Infinity and Choice are equiconsistent with TST with Infinity and TST with both Infinity and Choice, respectively. Therefore, a stronger theory like ZFC, which proves the consistency of TST, will also prove the consistency of NFU with these additions. In simpler terms, NFU is generally seen as weaker than NF because, in NFU, the collection of all sets (the power set of the universe) can be smaller than the universe itself, especially when urelements are included, as required by NFU with Choice.
Models of NFU
Jensen's proof gives a fairly simple method for producing models of NFU in bulk. Using well-known techniques of model theory, one can construct a nonstandard model of Zermelo set theory (nothing nearly as strong as full ZFC is needed for the basic technique) on which there is an external automorphism j (not a set of the model) which moves a rank of the cumulative hierarchy of sets. We may suppose without loss of generality that .
The domain of the model of NFU will be the nonstandard rank . The basic idea is that the automorphism j codes the "power set" of our "universe" into its externally isomorphic copy inside our "universe." The remaining objects not coding subsets of the universe are treated as urelements. Formally, the membership relation of the model of NFU will be
It may now be proved that this actually is a model of NFU. Let be a stratified formula in the language of NFU. Choose an assignment of types to all variables in the formula which witnesses the fact that it is stratified. Choose a natural number N greater than all types assigned to variables by this stratification. Expand the formula into a formula in the language of the nonstandard model of Zermelo set theory with automorphism j using the definition of membership in the model of NFU. Application of any power of j to both sides of an equation or membership statement preserves its truth value because j is an automorphism. Make such an application to each atomic formula in in such a way that each variable x assigned type i occurs with exactly applications of j. This is possible thanks to the form of the atomic membership statements derived from NFU membership statements, and to the formula being stratified. Each quantified sentence can be converted to the form (and similarly for existential quantifiers). Carry out this transformation everywhere and obtain a formula in which j is never applied to a bound variable. Choose any free variable y in assigned type i. Apply uniformly to the entire formula to obtain a formula in which y appears without any application of j. Now exists (because j appears applied only to free variables and constants), belongs to , and contains exactly those y which satisfy the original formula
in the model of NFU. has this extension in the model of NFU (the application of j corrects for the different definition of membership in the model of NFU). This establishes that Stratified Comprehension holds in the model of NFU.
To see that weak Extensionality holds is straightforward: each nonempty element of inherits a unique extension from the nonstandard model, the empty set inherits its usual extension as well, and all other objects are urelements.
If is a natural number n, one gets a model of NFU which claims that the universe is finite (it is externally infinite, of course). If is infinite and the Choice holds in the nonstandard model of ZFC, one obtains a model of NFU + Infinity + Choice.
Self-sufficiency of mathematical foundations in NFU
For philosophical reasons, it is important to note that it is not necessary to work in ZFC or any related system to carry out this proof. A common argument against the use of NFU as a foundation for mathematics is that the reasons for relying on it have to do with the intuition that ZFC is correct. It is sufficient to accept TST (in fact TSTU). In outline: take the type theory TSTU (allowing urelements in each positive type) as a metatheory and consider the theory of set models of TSTU in TSTU (these models will be sequences of sets (all of the same type in the metatheory) with embeddings of each into coding embeddings of the power set of into in a type-respecting manner). Given an embedding of into (identifying elements of the base "type" with subsets of the base type), embeddings may be defined from each "type" into its successor in a natural way. This can be generalized to transfinite sequences with care.
Note that the construction of such sequences of sets is limited by the size of the type in which they are being constructed; this prevents TSTU from proving its own consistency (TSTU + Infinity can prove the consistency of TSTU; to prove the consistency of TSTU+Infinity one needs a type containing a set of cardinality , which cannot be proved to exist in TSTU+Infinity without stronger assumptions). Now the same results of model theory can be used to build a model of NFU and verify that it is a model of NFU in much the same way, with the 's being used in place of in the usual construction. The final move is to observe that since NFU is consistent, we can drop the use of absolute types in our metatheory, bootstrapping the metatheory from TSTU to NFU.
Facts about the automorphism j
The automorphism j of a model of this kind is closely related to certain natural operations in NFU. For example, if W is a well-ordering in the nonstandard model (we suppose here that we use Kuratowski pairs so that the coding of functions in the two theories will agree to some extent) which is also a well-ordering in NFU (all well-orderings of NFU are well-orderings in the nonstandard model of Zermelo set theory, but not vice versa, due to the formation of urelements in the construction of the model), and W has type α in NFU, then j(W) will be a well-ordering of type T(α) in NFU.
In fact, j is coded by a function in the model of NFU. The function in the nonstandard model which sends the singleton of any element of to its sole element, becomes in NFU a function which sends each singleton {x}, where x is any object in the universe, to j(x). Call this function Endo and let it have the following properties: Endo is an injection from the set of singletons into the set of sets, with the property that Endo( {x} ) = {Endo( {y} ) | y∈x} for each set x. This function can define a type level "membership" relation on the universe, one reproducing the membership relation of the original nonstandard model.
History
In 1914, Norbert Wiener showed how to code the ordered pair as a set of sets, making it possible to eliminate the relation types of Principia Mathematica in favor of the linear hierarchy of sets in TST. The usual definition of the ordered pair was first proposed by Kuratowski in 1921. Willard Van Orman Quine first proposed NF as a way to avoid the "disagreeable consequences" of TST in a 1937 article titled New Foundations for Mathematical Logic; hence the name. Quine extended the theory in his book Mathematical Logic, whose first edition was published in 1940. In the book, Quine introduced the system "Mathematical Logic" or "ML", an extension of NF that included proper classes as well as sets. The first edition's set theory married NF to the proper classes of NBG set theory and included an axiom schema of unrestricted comprehension for proper classes. However, J. Barkley Rosser proved that the system was subject to the Burali-Forti paradox. Hao Wang showed how to amend Quine's axioms for ML so as to avoid this problem. Quine included the resulting axiomatization in the second and final edition, published in 1951.
In 1944, Theodore Hailperin showed that Comprehension is equivalent to a finite conjunction of its instances, In 1953, Ernst Specker showed that the axiom of choice is false in NF (without urelements). In 1969, Jensen showed that adding urelements to NF yields a theory (NFU) that is provably consistent. That same year, Grishin proved NF3 consistent. Specker additionally showed that NF is equiconsistent with TST plus the axiom scheme of "typical ambiguity". NF is also equiconsistent with TST augmented with a "type shifting automorphism", an operation (external to the theory) which raises type by one, mapping each type onto the next higher type, and preserves equality and membership relations.
In 1983, Marcel Crabbé proved consistent a system he called NFI, whose axioms are unrestricted extensionality and those instances of comprehension in which no variable is assigned a type higher than that of the set asserted to exist. This is a predicativity restriction, though NFI is not a predicative theory: it admits enough impredicativity to define the set of natural numbers (defined as the intersection of all inductive sets; note that the inductive sets quantified over are of the same type as the set of natural numbers being defined). Crabbé also discussed a subtheory of NFI, in which only parameters (free variables) are allowed to have the type of the set asserted to exist by an instance of comprehension. He called the result "predicative NF" (NFP); it is, of course, doubtful whether any theory with a self-membered universe is truly predicative. Holmes has shown that NFP has the same consistency strength as the predicative theory of types of Principia Mathematica without the axiom of reducibility.
The Metamath database implemented Hailperin's finite axiomatization for New Foundations. Since 2015, several candidate proofs by Randall Holmes of the consistency of NF relative to ZF were available both on arXiv and on the logician's home page. His proofs were based on demonstrating the equiconsistency of a "weird" variant of TST, "tangled type theory with λ-types" (TTTλ), with NF, and then showing that TTTλ is consistent relative to ZF with atoms but without choice (ZFA) by constructing a class model of ZFA which includes "tangled webs of cardinals" in ZF with atoms and choice (ZFA+C). These proofs were "difficult to read, insanely involved, and involve the sort of elaborate bookkeeping which makes it easy to introduce errors". In 2024, Sky Wilshaw formalized a version of Holmes' proof using the proof assistant Lean, finally resolving the question of NF's consistency. Timothy Chow characterized Wilshaw's work as showing that the reluctance of peer reviewers to engage with a difficult to understand proof can be addressed with the help of proof assistants.
| Mathematics | Axiomatic systems | null |
946319 | https://en.wikipedia.org/wiki/Orpiment | Orpiment | Orpiment, also known as ″yellow arsenic blende″ is a deep-colored, orange-yellow arsenic sulfide mineral with formula . It is found in volcanic fumaroles, low-temperature hydrothermal veins, and hot springs and may be formed through sublimation.
Orpiment takes its name from the Latin auripigmentum (aurum, "gold" + pigmentum, "pigment"), due to its deep-yellow color. Orpiment once was widely used in artworks, medicine, and other applications. Because of its toxicity and instability, its usage has declined.
Etymology
The Latin auripigmentum (aurum, "gold" + pigmentum, "pigment") referred both to its deep-yellow color and to the historical belief that it was thought to contain gold. The Latin term was used by Pliny in the first century after Christ.
The Greek for orpiment was arsenikon, deriving from the Greek word arsenikos, meaning "male", from the belief that metals were of different sexes. This Greek term was used by Theophrastus in the fourth century BC.
The Chinese term for orpiment is Ci-Huang (in Pinyin), meaning "female yellow".
The Persian for orpiment is zarnikh, deriving from the word "zar", the Persian for gold.
Physical and optical properties
Orpiment is a common monoclinic arsenic sulfide mineral. It has a Mohs hardness of 1.5 to 2 and a specific gravity of 3.49. It melts at to . Optically, it is biaxial (−) with refractive indices of a = 2.4, b = 2.81, g = 3.02.
Visual characteristics
Orpiment is a type of lemon-yellow to golden- or brownish-yellow crystal commonly found in foliated columnar or fibrous aggregates, may alternatively be botryoidal or reniform, granular or powdery, and, rarely, as prismatic crystals. Used as a pigment, orpiment's color is often described as a lemon- or canary-yellow, and occasionally as a golden- or brownish-yellow.
In the Munsell color system, "orpiment" is designated "brilliant yellow", Munsell notation 4.4Y 8.7 /8.9.
Orpiment and realgar
Orpiment and realgar are closely related minerals and are often categorized in the same group. They are both arsenic sulfides and belong to the monoclinic crystal system. They are found in the same deposits and can form in the same geologic environments. As a result, Orpiment and realgar share similar physical properties and histories of use by humans.
In Chinese, the names for orpiment and realgar are Ci-Huang and Xiong-Huang, respectively meaning "female yellow" and "male yellow". Their names symbolize their close natural conjunction, both physically in terms of their occurrence and properties and culturally in Chinese traditions.
Orpiment and realgar can be distinguished by their different visual characteristics. While orpiment typically has a vibrant golden-yellow color, realgar, in contrast, normally has an orange or reddish hue.
Permanence and conservation
Yellow orpiment (As2S3) degrades into arsenic oxides. Because of their solubility in water, arsenic oxides readily migrate to the surrounding environment. In painted works using orpiment, migrating, degraded arsenic oxides are often detectable throughout the multi-layered paint system. This widespread arsenic migration has consequences for the conservation of orpiment as a pigment in works of art.
Orpiment is also sensitive to light exposure, decaying into a friable white arsenic trioxide over time. Similarly, on ancient, orpiment-coated manuscript paper in Nepal, orpiment used to deter insects has often turned white over time.
Because of orpiment's solubility and instability as a pigment, preventing the degradation of orpiment may need to be prioritized in art conservation. Proper conservation methods should minimize exposure to strong light. Such methods should emphasize humidity control and avoid the use of water-based cleaning agents.
Use by artists
Orpiment has historically been used in artworks in many locales in the Eastern Hemisphere. It was one of the few clear, bright-yellow pigments available to artists until the 19th century.
Historical and regional use of orpiment
In Egypt, lumps of orpiment pigment have been found in a fourteenth-century BC tomb. In China, orpiment is known to have been used to color Chinese lacquer, despite no written sources mentioning this. Orpiment has also been identified on Central Asian wall paintings from the sixth to the thirteenth centuries. In a traditional Thai painting technique, still in use today, yellow ink for writing and drawing on black paper manuscripts is made using orpiment.
Medieval European artists imported orpiment from Asia Minor. Orpiment has been identified on Norwegian wooden altar frontals, polychrome sculptures, and folk art objects, including a crucifix. It was also used in twelfth- to sixteenth-century Eastern Orthodox icons from Bulgaria, Russia, and the former Yugoslavia. In Venice, records show that orpiment was purchased for a Romanian prince in 1600. European use of orpiment was uncommon until the nineteenth century, during which it saw use as a pigment in Impressionist paintings.
Orpiment as a pigment
In the Medieval Norwegian church of Tingelstad, orpiment was used in painting the altar frontal.
Orpiment was commonly combined with indigo dye to make a dark, rich green. In the Wilton Diptych (c 1395-9), this green pigment was used in egg tempera on the left panel for the green cloak of Edmund the Martyr.
Renaissance artists such as Raphael also used orpiment as a yellow pigment. In Raphael's Sistine Madonna from 1513–14, orpiment is used to achieve yellow on the clothing of the figures and in the background.
Tintoretto's Portrait of Vincenzo Morosini from about 1575–80 uses the pigment in its details. Orpiment is used to replicate the gold embroidery on Morosini's embroidered stole and to highlight the fur of the spotted ferret on his chest.
Limitations
Orpiment was one of the few clear, bright-yellow pigments available to artists until the 19th century. Its extreme toxicity and incompatibility with other, common, pigments, including lead and copper-based substances such as verdigris and azurite, meant that its use as a pigment ended when cadmium yellows, chromium yellows and organic aniline dye-based colors were introduced during the 19th century.
Other historical uses
Orpiment was traded in the Roman Empire and was used as a medicine in China, even though it is very toxic. It has been used as fly poison and to tip arrows with poison. Because of its striking color, it was of interest to alchemists, both in China and Europe, searching for a way to make gold. It also has been found in the wall decorations of Tutankhamun's tomb and ancient Egyptian scrolls, and on the walls of the Taj Mahal. For centuries, orpiment was ground down and used as a pigment in painting and for sealing wax, and was even used in ancient China as a correction fluid. Orpiment is mentioned in the 17th century by Robert Hooke in Micrographia for the manufacture of small shot. Scientists like Richard Adolf Zsigmondy and Hermann Ambronn puzzled jointly over the amorphous form of , "orpiment glass", as early as 1904.
Industry uses
Orpiment is used in the production of infrared-transmitting glass, oil cloth, linoleum, semiconductors, photoconductors, pigments, and fireworks. Mixed with two parts of slaked lime (calcium hydroxide), orpiment is still commonly used in rural India as a depilatory. It is used in the tanning industry to remove hair from hides.
Orpiment has been used as bookends. In 2023, the UK Office for Product Safety and Standards recalled 40 pieces sold by TK Maxx between June and October 2022, due to the mineral's toxicity.
Crystal structure
Gallery
| Physical sciences | Minerals | Earth science |
946975 | https://en.wikipedia.org/wiki/Line%20%28geometry%29 | Line (geometry) | In geometry, a straight line, usually abbreviated line, is an infinitely long object with no width, depth, or curvature, an idealization of such physical objects as a straightedge, a taut string, or a ray of light. Lines are spaces of dimension one, which may be embedded in spaces of dimension two, three, or higher. The word line may also refer, in everyday life, to a line segment, which is a part of a line delimited by two points (its endpoints).
Euclid's Elements defines a straight line as a "breadthless length" that "lies evenly with respect to the points on itself", and introduced several postulates as basic unprovable properties on which the rest of geometry was established. Euclidean line and Euclidean geometry are terms introduced to avoid confusion with generalizations introduced since the end of the 19th century, such as non-Euclidean, projective, and affine geometry.
Properties
In the Greek deductive geometry of Euclid's Elements, a general line (now called a curve) is defined as a "breadthless length", and a straight line (now called a line segment) was defined as a line "which lies evenly with the points on itself". These definitions appeal to readers' physical experience, relying on terms that are not themselves defined, and the definitions are never explicitly referenced in the remainder of the text. In modern geometry, a line is usually either taken as a primitive notion with properties given by axioms, or else defined as a set of points obeying a linear relationship, for instance when real numbers are taken to be primitive and geometry is established analytically in terms of numerical coordinates.
In an axiomatic formulation of Euclidean geometry, such as that of Hilbert (modern mathematicians added to Euclid's original axioms to fill perceived logical gaps), a line is stated to have certain properties that relate it to other lines and points. For example, for any two distinct points, there is a unique line containing them, and any two distinct lines intersect at most at one point. In two dimensions (i.e., the Euclidean plane), two lines that do not intersect are called parallel. In higher dimensions, two lines that do not intersect are parallel if they are contained in a plane, or skew if they are not.
On a Euclidean plane, a line can be represented as a boundary between two regions. Any collection of finitely many lines partitions the plane into convex polygons (possibly unbounded); this partition is known as an arrangement of lines.
In higher dimensions
In three-dimensional space, a first degree equation in the variables x, y, and z defines a plane, so two such equations, provided the planes they give rise to are not parallel, define a line which is the intersection of the planes. More generally, in n-dimensional space n−1 first-degree equations in the n coordinate variables define a line under suitable conditions.
In more general Euclidean space, Rn (and analogously in every other affine space), the line L passing through two different points a and b is the subset
The direction of the line is from a reference point a (t = 0) to another point b (t = 1), or in other words, in the direction of the vector b − a. Different choices of a and b can yield the same line.
Collinear points
Three or more points are said to be collinear if they lie on the same line. If three points are not collinear, there is exactly one plane that contains them.
In affine coordinates, in n-dimensional space the points X = (x1, x2, ..., xn), Y = (y1, y2, ..., yn), and Z = (z1, z2, ..., zn) are collinear if the matrix
has a rank less than 3. In particular, for three points in the plane (n = 2), the above matrix is square and the points are collinear if and only if its determinant is zero.
Equivalently for three points in a plane, the points are collinear if and only if the slope between one pair of points equals the slope between any other pair of points (in which case the slope between the remaining pair of points will equal the other slopes). By extension, k points in a plane are collinear if and only if any (k–1) pairs of points have the same pairwise slopes.
In Euclidean geometry, the Euclidean distance d(a,b) between two points a and b may be used to express the collinearity between three points by:
The points a, b and c are collinear if and only if d(x,a) = d(c,a) and d(x,b) = d(c,b) implies x = c.
However, there are other notions of distance (such as the Manhattan distance) for which this property is not true.
In the geometries where the concept of a line is a primitive notion, as may be the case in some synthetic geometries, other methods of determining collinearity are needed.
Relationship with other figures
In Euclidean geometry, all lines are congruent, meaning that every line can be obtained by moving a specific line. However, lines may play special roles with respect to other geometric objects and can be classified according to that relationship.
For instance, with respect to a conic (a circle, ellipse, parabola, or hyperbola), lines can be:
tangent lines, which touch the conic at a single point;
secant lines, which intersect the conic at two points and pass through its interior;
exterior lines, which do not meet the conic at any point of the Euclidean plane; or
a directrix, whose distance from a point helps to establish whether the point is on the conic.
a coordinate line, a linear coordinate dimension
In the context of determining parallelism in Euclidean geometry, a transversal is a line that intersects two other lines that may or not be parallel to each other.
For more general algebraic curves, lines could also be:
i-secant lines, meeting the curve in i points counted without multiplicity, or
asymptotes, which a curve approaches arbitrarily closely without touching it.
With respect to triangles we have:
the Euler line,
the Simson lines, and
central lines.
For a convex quadrilateral with at most two parallel sides, the Newton line is the line that connects the midpoints of the two diagonals.
For a hexagon with vertices lying on a conic we have the Pascal line and, in the special case where the conic is a pair of lines, we have the Pappus line.
Parallel lines are lines in the same plane that never cross. Intersecting lines share a single point in common. Coincidental lines coincide with each other—every point that is on either one of them is also on the other.
Perpendicular lines are lines that intersect at right angles.
In three-dimensional space, skew lines are lines that are not in the same plane and thus do not intersect each other.
In axiomatic systems
In synthetic geometry, the concept of a line is often considered as a primitive notion, meaning it is not being defined by using other concepts, but it is defined by the properties, called axioms, that it must satisfy.
However, the axiomatic definition of a line does not explain the relevance of the concept and is often too abstract for beginners. So, the definition is often replaced or completed by a mental image or intuitive description that allows understanding what is a line. Such descriptions are sometimes referred to as definitions, but are not true definitions since they cannot used in mathematical proofs. The "definition" of line in Euclid's Elements falls into this category; and is never used in proofs of theorems.
Definition
Linear equation
Lines in a Cartesian plane or, more generally, in affine coordinates, are characterized by linear equations. More precisely, every line (including vertical lines) is the set of all points whose coordinates (x, y) satisfy a linear equation; that is,
where a, b and c are fixed real numbers (called coefficients) such that a and b are not both zero. Using this form, vertical lines correspond to equations with b = 0.
One can further suppose either or , by dividing everything by if it is not zero.
There are many variant ways to write the equation of a line which can all be converted from one to another by algebraic manipulation. The above form is sometimes called the standard form. If the constant term is put on the left, the equation becomes
and this is sometimes called the general form of the equation. However, this terminology is not universally accepted, and many authors do not distinguish these two forms.
These forms are generally named by the type of information (data) about the line that is needed to write down the form. Some of the important data of a line is its slope, x-intercept, known points on the line and y-intercept.
The equation of the line passing through two different points and may be written as
If , this equation may be rewritten as
or
In two dimensions, the equation for non-vertical lines is often given in the slope–intercept form:
where:
m is the slope or gradient of the line.
b is the y-intercept of the line.
x is the independent variable of the function .
The slope of the line through points and , when , is given by and the equation of this line can be written .
As a note, lines in three dimensions may also be described as the simultaneous solutions of two linear equations
such that and are not proportional (the relations imply ). This follows since in three dimensions a single linear equation typically describes a plane and a line is what is common to two distinct intersecting planes.
Parametric equation
Parametric equations are also used to specify lines, particularly in those in three dimensions or more because in more than two dimensions lines cannot be described by a single linear equation.
In three dimensions lines are frequently described by parametric equations:
where:
x, y, and z are all functions of the independent variable t which ranges over the real numbers.
(x0, y0, z0) is any point on the line.
a, b, and c are related to the slope of the line, such that the direction vector (a, b, c) is parallel to the line.
Parametric equations for lines in higher dimensions are similar in that they are based on the specification of one point on the line and a direction vector.
Hesse normal form
The normal form (also called the Hesse normal form, after the German mathematician Ludwig Otto Hesse), is based on the normal segment for a given line, which is defined to be the line segment drawn from the origin perpendicular to the line. This segment joins the origin with the closest point on the line to the origin. The normal form of the equation of a straight line on the plane is given by:
where is the angle of inclination of the normal segment (the oriented angle from the unit vector of the -axis to this segment), and is the (positive) length of the normal segment. The normal form can be derived from the standard form by dividing all of the coefficients by
and also multiplying through by if
Unlike the slope-intercept and intercept forms, this form can represent any line but also requires only two finite parameters, and , to be specified. If , then is uniquely defined modulo . On the other hand, if the line is through the origin (), one drops the term to compute and , and it follows that is only defined modulo .
Other representations
Vectors
The vector equation of the line through points A and B is given by (where λ is a scalar).
If a is vector OA and b is vector OB, then the equation of the line can be written: .
A ray starting at point A is described by limiting λ. One ray is obtained if λ ≥ 0, and the opposite ray comes from λ ≤ 0.
Polar coordinates
In a Cartesian plane, polar coordinates are related to Cartesian coordinates by the parametric equations:
In polar coordinates, the equation of a line not passing through the origin—the point with coordinates —can be written
with and Here, is the (positive) length of the line segment perpendicular to the line and delimited by the origin and the line, and is the (oriented) angle from the -axis to this segment.
It may be useful to express the equation in terms of the angle between the -axis and the line. In this case, the equation becomes
with and
These equations can be derived from the normal form of the line equation by setting and and then applying the angle difference identity for sine or cosine.
These equations can also be proven geometrically by applying right triangle definitions of sine and cosine to the right triangle that has a point of the line and the origin as vertices, and the line and its perpendicular through the origin as sides.
The previous forms do not apply for a line passing through the origin, but a simpler formula can be written: the polar coordinates of the points of a line passing through the origin and making an angle of with the -axis, are the pairs such that
Generalizations
In modern mathematics, given the multitude of geometries, the concept of a line is closely tied to the way the geometry is described. For instance, in analytic geometry, a line in the plane is often defined as the set of points whose coordinates satisfy a given linear equation, but in a more abstract setting, such as incidence geometry, a line may be an independent object, distinct from the set of points which lie on it.
When a geometry is described by a set of axioms, the notion of a line is usually left undefined (a so-called primitive object). The properties of lines are then determined by the axioms which refer to them. One advantage to this approach is the flexibility it gives to users of the geometry. Thus in differential geometry, a line may be interpreted as a geodesic (shortest path between points), while in some projective geometries, a line is a 2-dimensional vector space (all linear combinations of two independent vectors). This flexibility also extends beyond mathematics and, for example, permits physicists to think of the path of a light ray as being a line.
Projective geometry
In many models of projective geometry, the representation of a line rarely conforms to the notion of the "straight curve" as it is visualised in Euclidean geometry. In elliptic geometry we see a typical example of this. In the spherical representation of elliptic geometry, lines are represented by great circles of a sphere with diametrically opposite points identified. In a different model of elliptic geometry, lines are represented by Euclidean planes passing through the origin. Even though these representations are visually distinct, they satisfy all the properties (such as, two points determining a unique line) that make them suitable representations for lines in this geometry.
The "shortness" and "straightness" of a line, interpreted as the property that the distance along the line between any two of its points is minimized (see triangle inequality), can be generalized and leads to the concept of geodesics in metric spaces.
Related concepts
Ray
Given a line and any point A on it, we may consider A as decomposing this line into two parts.
Each such part is called a ray and the point A is called its initial point. It is also known as half-line (sometimes, a half-axis if it plays a distinct role, e.g., as part of a coordinate axis). It is a one-dimensional half-space. The point A is considered to be a member of the ray. Intuitively, a ray consists of those points on a line passing through A and proceeding indefinitely, starting at A, in one direction only along the line. However, in order to use this concept of a ray in proofs a more precise definition is required.
Given distinct points A and B, they determine a unique ray with initial point A. As two points define a unique line, this ray consists of all the points between A and B (including A and B) and all the points C on the line through A and B such that B is between A and C. This is, at times, also expressed as the set of all points C on the line determined by A and B such that A is not between B and C. A point D, on the line determined by A and B but not in the ray with initial point A determined by B, will determine another ray with initial point A. With respect to the AB ray, the AD ray is called the opposite ray.
Thus, we would say that two different points, A and B, define a line and a decomposition of this line into the disjoint union of an open segment and two rays, BC and AD (the point D is not drawn in the diagram, but is to the left of A on the line AB). These are not opposite rays since they have different initial points.
In Euclidean geometry two rays with a common endpoint form an angle.
The definition of a ray depends upon the notion of betweenness for points on a line. It follows that rays exist only for geometries for which this notion exists, typically Euclidean geometry or affine geometry over an ordered field. On the other hand, rays do not exist in projective geometry nor in a geometry over a non-ordered field, like the complex numbers or any finite field.
Line segment
A line segment is a part of a line that is bounded by two distinct end points and contains every point on the line between its end points. Depending on how the line segment is defined, either of the two end points may or may not be part of the line segment. Two or more line segments may have some of the same relationships as lines, such as being parallel, intersecting, or skew, but unlike lines they may be none of these, if they are coplanar and either do not intersect or are collinear.
Number line
A point on number line corresponds to a real number and vice versa. Usually, integers are evenly spaced on the line, with positive numbers are on the right, negative numbers on the left. As an extension to the concept, an imaginary line representing imaginary numbers can be drawn perpendicular to the number line at zero. The two lines forms the complex plane, a geometrical representation of the set of complex numbers.
| Mathematics | Geometry and topology | null |
947038 | https://en.wikipedia.org/wiki/Bird%20nest | Bird nest | A bird nest is the spot in which a bird lays and incubates its eggs and raises its young. Although the term popularly refers to a specific structure made by the bird itself—such as the grassy cup nest of the American robin or Eurasian blackbird, or the elaborately woven hanging nest of the Montezuma oropendola or the village weaver—that is too restrictive a definition. For some species, a nest is simply a shallow depression made in sand; for others, it is the knot-hole left by a broken branch, a burrow dug into the ground, a chamber drilled into a tree, an enormous rotting pile of vegetation and earth, a shelf made of dried saliva or a mud dome with an entrance tunnel. The smallest bird nests are those of some hummingbirds, tiny cups which can be a mere across and high. At the other extreme, some nest mounds built by the dusky scrubfowl measure more than in diameter and stand nearly tall. The study of birds' nests is known as caliology.
Not all bird species build nests. Some species lay their eggs directly on the ground or rocky ledges, while brood parasites lay theirs in the nests of other birds, letting unwitting "foster parents" do the work of rearing the young. Although nests are primarily used for breeding, they may also be reused in the non-breeding season for roosting and some species build special dormitory nests or roost nests (or winter-nest) that are used only for roosting. Most birds build a new nest each year, though some refurbish their old nests. The large eyries (or aeries) of some eagles are platform nests that have been used and refurbished for several years.
In the majority of nest-building species the female does most or all of the nest construction, in others both partners contribute; sometimes the male builds the nest and the hen lines it. In some polygynous species, however, the male does most or all of the nest building. The nest may also form a part of the courtship display such as in weaver birds. The ability to choose and maintain good nest sites and build high quality nests may be selected for by females in these species. In some species the young from previous broods may also act as helpers for the adults.
Type
Not every bird species builds or uses a nest. Some auks, for instance—including common murre, thick-billed murre and razorbill—lay their eggs directly onto the narrow rocky ledges they use as breeding sites. The eggs of these species are dramatically pointed at one end, so that they roll in a circle when disturbed. This is critical for the survival of the developing eggs, as there are no nests to keep them from rolling off the side of the cliff. Presumably because of the vulnerability of their unprotected eggs, parent birds of these auk species rarely leave them unattended. Nest location and architecture is strongly influenced by local topography and other abiotic factors.
King penguins and emperor penguins also do not build nests; instead, they tuck their eggs and chicks between their feet and folds of skin on their lower bellies. They are thus able to move about while incubating, though in practice only the emperor penguin regularly does so. Emperor penguins breed during the harshest months of the Antarctic winter, and their mobility allows them to form huge huddled masses which help them to withstand the extremely high winds and low temperatures of the season. Without the ability to share body heat (temperatures in the centre of tight groups can be as much as 10C above the ambient air temperature), the penguins would expend far more energy trying to stay warm, and breeding attempts would probably fail.
Some crevice-nesting species, including ashy storm-petrel, pigeon guillemot, Eurasian eagle-owl and Hume's tawny owl, lay their eggs in the relative shelter of a crevice in the rocks or a gap between boulders, but provide no additional nest material. Potoos lay their single egg directly atop a broken stump, or into a shallow depression on a branch—typically where an upward-pointing branch died and fell off, leaving a small scar or knot-hole. Brood parasites, such as the New World cowbirds, the honeyguides, and many of the Old World and Australasian cuckoos, lay their eggs in the active nests of other species.
Scrape
The simplest nest construction is the scrape, which is merely a shallow depression in soil or vegetation. This nest type, which typically has a rim deep enough to keep the eggs from rolling away, is sometimes lined with bits of vegetation, small stones, shell fragments or feathers. These materials may help to camouflage the eggs or may provide some level of insulation; they may also help to keep the eggs in place, and prevent them from sinking into muddy or sandy soil if the nest is accidentally flooded. Ostriches, most tinamous, many ducks, most shorebirds, most terns, some falcons, pheasants, quail, partridges, bustards and sandgrouse are among the species that build scrape nests.
Eggs and young in scrape nests, and the adults that brood them, are more exposed to predators and the elements than those in more sheltered nests; they are on the ground and typically in the open, with little to hide them. The eggs of most ground-nesting birds (including those that use scrape nests) are cryptically coloured to help camouflage them when the adult is not covering them; the actual colour generally corresponds to the substrate on which they are laid. Brooding adults also tend to be well camouflaged, and may be difficult to flush from the nest. Most ground-nesting species have well-developed distraction displays, which are used to draw (or drive) potential predators from the area around the nest. Most species with this type of nest have precocial young, which quickly leave the nest upon hatching.
In cool climates (such as in the high Arctic or at high elevations), the depth of a scrape nest can be critical to both the survival of developing eggs and the fitness of the parent bird incubating them. The scrape must be deep enough that eggs are protected from the convective cooling caused by cold winds, but shallow enough that they and the parent bird are not too exposed to the cooling influences of ground temperatures, particularly where the permafrost layer rises to mere centimeters below the nest. Studies have shown that an egg within a scrape nest loses heat 9% more slowly than an egg placed on the ground beside the nest; in such a nest lined with natural vegetation, heat loss is reduced by an additional 25%. The insulating factor of nest lining is apparently so critical to egg survival that some species, including Kentish plovers, will restore experimentally altered levels of insulation to their pre-adjustment levels (adding or subtracting material as necessary) within 24 hours.
In warm climates, such as deserts and salt flats, heat rather than cold can kill the developing embryos. In such places, scrapes are shallower and tend to be lined with non-vegetative material (including shells, feathers, sticks and soil), which allows convective cooling to occur as air moves over the eggs. Some species, such as the lesser nighthawk and the red-tailed tropicbird, help reduce the nest's temperature by placing it in partial or full shade. Others, including some shorebirds, cast shade with their bodies as they stand over their eggs. Some shorebirds also soak their breast feathers with water and then sit on the eggs, providing moisture to enable evaporative cooling. Parent birds keep from overheating themselves by gular panting while they are incubating, frequently exchanging incubation duties, and standing in water when they are not incubating.
The technique used to construct a scrape nest varies slightly depending on the species. Beach-nesting terns, for instance, fashion their nests by rocking their bodies on the sand in the place they have chosen to site their nest, while skimmers build their scrapes with their feet, kicking sand backwards while resting on their bellies and turning slowly in circles. The ostrich also scratches out its scrape with its feet, though it stands while doing so. Many tinamous lay their eggs on a shallow mat of dead leaves they have collected and placed under bushes or between the root buttresses of trees, and kagus lay theirs on a pile of dead leaves against a log, tree trunk or vegetation. Marbled godwits stomp a grassy area flat with their feet, then lay their eggs, while other grass-nesting waders bend vegetation over their nests so as to avoid detection from above. Many female ducks, particularly in the northern latitudes, line their shallow scrape nests with down feathers plucked from their own breasts, as well as with small amounts of vegetation. Among scrape-nesting birds, the three-banded courser and Egyptian plover are unique in their habit of partially burying their eggs in the sand of their scrapes.
Mound
Burying eggs as a form of incubation reaches its zenith with the Australasian megapodes. Several megapode species construct enormous mound nests made of soil, branches, sticks, twigs and leaves, and lay their eggs within the rotting mass. The heat generated by these mounds, which are in effect giant compost heaps, warms and incubates the eggs. The nest heat results from the respiration of thermophilic fungi and other microorganisms. The size of some of these mounds can be truly staggering; several of the largest—which contain more than of material, and probably weigh more than 50 tons (45,000 kg)—were initially thought to be Aboriginal middens.
In most mound-building species, males do most or all of the nest construction and maintenance. Using his strong legs and feet, the male scrapes together material from the area around his chosen nest site, gradually building a conical or bell-shaped pile. This process can take five to seven hours a day for more than a month. While mounds are typically reused for multiple breeding seasons, new material must be added each year to generate the appropriate amount of heat. A female will begin to lay eggs in the nest only when the mound's temperature has reached an optimal level.
Both the temperature and the moisture content of the mound are critical to the survival and development of the eggs, so both are carefully regulated for the entire length of the breeding season (which may last for as long as eight months), principally by the male. Ornithologists believe that megapodes may use sensitive areas in their mouths to assess mound temperatures; each day during the breeding season, the male digs a pit into his mound and sticks his head in. If the mound's core temperature is a bit low, he adds fresh moist material to the mound, and stirs it in; if it is too high, he opens the top of the mound to allow some of the excess heat to escape. This regular monitoring also keeps the mound's material from becoming compacted, which would inhibit oxygen diffusion to the eggs and make it more difficult for the chicks to emerge after hatching. The malleefowl, which lives in more open forest than do other megapodes, uses the sun to help warm its nest as well—opening the mound at midday during the cool spring and autumn months to expose the plentiful sand incorporated into the nest to the sun's warming rays, then using that warm sand to insulate the eggs during the cold nights. During hot summer months, the malleefowl opens its nest mound only in the cool early morning hours, allowing excess heat to escape before recovering the mound completely. One recent study showed that the sex ratio of Australian brushturkey hatchlings correlated strongly with mound temperatures; females hatched from eggs incubated at higher mean temperatures.
Flamingos make a different type of mound nest. Using their beaks to pull material towards them, they fashion a cone-shaped pile of mud between tall, with a small depression in the top to house their single egg. The height of the nest varies with the substrate upon which it is built; those on clay sites are taller on average than those on dry or sandy sites. The height of the nest and the circular, often water-filled trench which surrounds it (the result of the removal of material for the nest) help to protect the egg from fluctuating water levels and excessive heat at ground level. In East Africa, for example, temperatures at the top of the nest mound average some cooler than those of the surrounding ground.
The base of the horned coot's enormous nest is a mound built of stones, gathered one at a time by the pair, using their beaks. These stones, which may weigh as much as 450 g (about a pound) each, are dropped into the shallow water of a lake, making a cone-shaped pile which can measure as much as at the bottom and at the top, and in height. The total combined weight of the mound's stones may approach 1.5 tons (1,400 kg). Once the mound has been completed, a sizable platform of aquatic vegetation is constructed on top. The entire structure is typically reused for many years.
Burrow
Soil plays a different role in the burrow nest; here, the eggs and young—and in most cases the incubating parent bird—are sheltered under the earth. Most burrow-nesting birds excavate their own burrows, but some use those excavated by other species and are known as secondary nesters; burrowing owls, for example, sometimes use the burrows of prairie dogs, ground squirrels, badgers or tortoises, China's endemic white-browed tits use the holes of ground-nesting rodents and common kingfishers occasionally nest in rabbit burrows. Burrow nests are particularly common among seabirds at high latitudes, as they provide protection against both cold temperatures and predators. Puffins, shearwaters, some megapodes, motmots, todies, most kingfishers, the crab plover, miners and leaftossers are among the species which use burrow nests.
Most burrow nesting species dig a horizontal tunnel into a vertical (or nearly vertical) dirt cliff, with a chamber at the tunnel's end to house the eggs. The length of the tunnel varies depending on the substrate and the species; sand martins make relatively short tunnels ranging from , for example, while those of the burrowing parakeet can extend for more than three meters (nearly 10 ft). Some species, including the ground-nesting puffbirds, prefer flat or gently sloping land, digging their entrance tunnels into the ground at an angle. In a more extreme example, the D'Arnaud's barbet digs a vertical tunnel shaft more than a meter (39 in) deep, with its nest chamber excavated off to the side at some height above the shaft's bottom; this arrangement helps to keep the nest from being flooded during heavy rain. Buff-breasted paradise-kingfishers dig their nests into the compacted mud of active termite mounds, either on the ground or in trees. Specific soil types may favour certain species and it is speculated that several species of bee-eater favor loess soils which are easy to penetrate.
Birds use a combination of their beaks and feet to excavate burrow nests. The tunnel is started with the beak; the bird either probes at the ground to create a depression, or flies toward its chosen nest site on a cliff wall and hits it with its bill. The latter method is not without its dangers; there are reports of kingfishers being fatally injured in such attempts. Some birds remove tunnel material with their bills, while others use their bodies or shovel the dirt out with one or both feet. Female paradise-kingfishers are known to use their long tails to clear the loose soil.
Some crepuscular petrels and prions are able to identify their own burrows within dense colonies by smell. Sand martins learn the location of their nest within a colony, and will accept any chick put into that nest until right before the young fledge.
Not all burrow-nesting species incubate their young directly. Some megapode species bury their eggs in sandy pits dug where sunlight, subterranean volcanic activity, or decaying tree roots will warm the eggs. The crab plover also uses a burrow nest, the warmth of which allows it to leave the eggs unattended for as long as 58 hours.
Predation levels on some burrow-nesting species can be quite high; on Alaska's Wooded Islands, for example, river otters munched their way through some 23 percent of the island's fork-tailed storm-petrel population during a single breeding season in 1977. There is some evidence that increased vulnerability may lead some burrow-nesting species to form colonies, or to nest closer to rival pairs in areas of high predation than they might otherwise do.
Cavity
The cavity nest is a chamber, typically in living or dead wood, but sometimes in the trunks of tree ferns or large cacti, including saguaro. In tropical areas, cavities are sometimes excavated in arboreal insect nests. A relatively small number of species, including woodpeckers, trogons, some nuthatches and many barbets, can excavate their own cavities. Far more species—including parrots, tits, bluebirds, most hornbills, some kingfishers, some owls, some ducks and some flycatchers—use natural cavities, or those abandoned by species able to excavate them; they also sometimes usurp cavity nests from their excavating owners. Those species that excavate their own cavities are known as "primary cavity nesters", while those that use natural cavities or those excavated by other species are called "secondary cavity nesters". Both primary and secondary cavity nesters can be enticed to use nest boxes (also known as bird houses); these mimic natural cavities, and can be critical to the survival of species in areas where natural cavities are lacking.
Woodpeckers use their chisel-like bills to excavate their cavity nests, a process which takes, on average, about two weeks. Cavities are normally excavated on the downward-facing side of a branch, presumably to make it more difficult for predators to access the nest, and to reduce the chance that rain floods the nest. There is also some evidence that fungal rot may make the wood on the underside of leaning trunks and branches easier to excavate. Most woodpeckers use a cavity for only a single year. The endangered red-cockaded woodpecker is an exception; it takes far longer—up to two years—to excavate its nest cavity, and may reuse it for more than two decades. The typical woodpecker nest has a short horizontal tunnel which leads to a vertical chamber within the trunk. The size and shape of the chamber depends on species, and the entrance hole is typically only as large as is needed to allow access for the adult birds. While wood chips are removed during the excavation process, most species line the floor of the cavity with a fresh bed of them before laying their eggs.
Trogons excavate their nests by chewing cavities into very soft dead wood; some species make completely enclosed chambers (accessed by upward-slanting entrance tunnels), while others—like the extravagantly plumed resplendent quetzal—construct more open niches. In most trogon species, both sexes help with nest construction. The process may take several months, and a single pair may start several excavations before finding a tree or stump with wood of the right consistency.
Species which use natural cavities or old woodpecker nests sometimes line the cavity with soft material such as grass, moss, lichen, feathers or fur. Though a number of studies have attempted to determine whether secondary cavity nesters preferentially choose cavities with entrance holes facing certain directions, the results remain inconclusive. While some species appear to preferentially choose holes with certain orientations, studies (to date) have not shown consistent differences in fledging rates between nests oriented in different directions.
Cavity-dwelling species have to contend with the danger of predators accessing their nest, catching them and their young inside and unable to get out. They have a variety of methods for decreasing the likelihood of this happening. Red-cockaded woodpeckers peel bark around the entrance, and drill wells above and below the hole; since they nest in live trees, the resulting flow of resin forms a barrier that prevents snakes from reaching the nests. Red-breasted nuthatches smear sap around the entrance holes to their nests, while white-breasted nuthatches rub foul-smelling insects around theirs. Eurasian nuthatches wall up part of their entrance holes with mud, decreasing the size and sometimes extending the tunnel part of the chamber. Most female hornbills seal themselves into their cavity nests, using a combination of mud (in some species brought by their mates), food remains and their own droppings to reduce the entrance hole to a narrow slit.
A few birds are known to use the nests of insects within which they create a cavity in which they lay their eggs. These include the rufous woodpecker which nests in the arboreal nests of Crematogaster ants and the collared kingfisher which uses termite nests.
Cup
The cup nest is smoothly hemispherical inside, with a deep depression to house the eggs. Most are made of pliable materials—including grasses—though a small number are made of mud or saliva. Many passerines and a few non-passerines, including some hummingbirds and some swifts, build this type of nest.
Small bird species in more than 20 passerine families, and a few non-passerines—including most hummingbirds, kinglets and crests in the genus Regulus, some tyrant flycatchers and several New World warblers—use considerable amounts of spider silk in the construction of their nests. The lightweight material is strong and extremely flexible, allowing the nest to mold to the adult during incubation (reducing heat loss), then to stretch to accommodate the growing nestlings; as it is sticky, it also helps to bind the nest to the branch or leaf to which it is attached.
Many swifts and some hummingbirds use thick, quick-drying saliva to anchor their nests. The chimney swift starts by dabbing two globs of saliva onto the wall of a chimney or tree trunk. In flight, it breaks a small twig from a tree and presses it into the saliva, angling the twig downwards so that the central part of the nest is the lowest. It continues adding globs of saliva and twigs until it has made a crescent-shaped cup.
Cup-shaped nest insulation has been found to be related to nest mass, nest wall thickness, nest depth, nest weave density/porosity, surface area, height above ground and elevation above sea level.
More recently, nest insulation has been found to be related to the mass of the incubating parent. This is known as an allometric relationship. Nest walls are constructed with an adequate quantity of nesting material so that the nest will be capable of supporting the contents of the nest. Nest thickness, nest mass and nest dimensions therefore correlate with the mass of the adult bird. The flow-on consequence of this is that nest insulation is also related to parent mass.
Saucer or plate
The saucer or plate nest, though superficially similar to a cup nest, has at most only a shallow depression to house the eggs.
Platform
The platform nest is a large structure, often many times the size of the (typically large) bird which has built it. Depending on the species, these nests can be on the ground or elevated.
In the case of raptor nests, or eyries (also spelled aerie), these are often used for many years, with new material added each breeding season. In some cases, the nests grow large enough to cause structural damage to the tree itself, particularly during bad storms where the weight of the nest can cause additional stress on wind-tossed branches.
Pendent
The pendent nest is an elongated sac woven of pliable materials such as grasses and plant fibers and suspended from a branch. Oropendolas, caciques, orioles, weavers and sunbirds are among the species that weave pendent nests. In weaver birds, this is pendant, suspended from a single point hanging from branch while many other birds incorporate more than one branch to support the nest.
Sphere
The sphere nest is a roundish structure; it is completely enclosed, except for a small opening which allows access. Most spherical nests are woven out of plant material. Spider webs are also frequently used, upon which other material such as lichens may be stuck for camouflage. The cape penduline tit incorporates false entrances, the parent bird carefully making sure to close the actual entrance when leaving the nest. The entrances are lined with spider webs which help seal the openings.
Nest protection and sanitation
Many species of bird conceal their nests to protect them from predators. Some species may choose nest sites that are inaccessible or build the nest so as to deter predators. Bird nests can also act as habitats for other inquiline species which may not affect the bird directly. Birds have also evolved nest sanitation measures to reduce the effects of parasites and pathogens on nestlings.
Some aquatic species such as grebes are very careful when approaching and leaving the nest so as not to reveal the location. Some species will use leaves to cover up the nest prior to leaving.
Ground birds such as plovers may use broken wing or rodent run displays to distract predators from nests.
Many species attack predators or apparent predators near their nests. Kingbirds attack other birds that come too close. In North America, northern mockingbirds, blue jays, and Arctic terns can peck hard enough to draw blood. In Australia, a bird attacking a person near its nest is said to swoop the person. The Australian magpie is particularly well known for this behavior.
Nests can become home to many other organisms including parasites and pathogens. The excreta of the fledglings also pose a problem. In most passerines, the adults actively dispose the fecal sacs of young at a distance or consume them. This is believed to help prevent ground predators from detecting nests. Young birds of prey however usually void their excreta beyond the rims of their nests. Blowflies of the genus Protocalliphora have specialized to become obligate nest parasites with the maggots feeding on the blood of nestlings.
Some birds have been shown to choose aromatic green plant material for constructing nests that may have insecticidal properties, while others may use materials such as carnivore scat to repel smaller predators. Some urban birds, house sparrows and house finches in Mexico, have adopted the use of cigarette butts which contain nicotine and other toxic substances that repel ticks and other ectoparasites.
Some birds use pieces of snake slough in their nests. It has been suggested that these may deter some nest predators such as squirrels.
Colonial nesting
Though most birds nest individually, some species—including seabirds, penguins, flamingos, many herons, gulls, terns, weaver, some corvids and some sparrows—gather together in sizeable colonies. Birds that nest colonially may benefit from increased protection against predation. They may also be able to better use food supplies, by following more successful foragers to their foraging sites.
Ecological importance
In constructing nests, birds act as ecosystem engineers by providing a sheltered microclimate and concentrated food sources for invertebrates. A global checklist lists eighteen invertebrate orders that occur in bird nests.
In human culture
Many birds may nest close to human habitations. In addition to nest boxes which are often used to encourage cavity nesting birds (see below), other species have been specially encouraged : for example nesting white storks have been protected and held in reverence in many cultures, and the nesting of peregrine falcons on tall modern or historical buildings has captured popular interest.
Colonial breeders produce guano in and around their nesting sites, which is a valuable fertilizer from the Andean Pacific coast and other areas.
The saliva nest of the edible-nest swiftlet is used to make bird's nest soup, long considered a delicacy in China. Collection of the swiftlet nests is big business: in one year, more than 3.5 million nests were exported from Borneo to China, and the industry was estimated at $1 billion US per year (and increasing) in 2008. While the collection is regulated in some areas (at the Gomantong Caves, for example, where nests can be collected only from February to April or July to September), it is not in others, and the swiftlets are declining in areas where the harvest reaches unsustainable levels.
Some species of birds are considered nuisances when they nest in the proximity of human habitations. Feral pigeons are often unwelcome and sometimes also considered as a health risk.
The Beijing National Stadium, principal venue of the 2008 Summer Olympics, has been nicknamed "The Bird Nest" because of its architectural design, which its designers likened to a bird's woven nest.
In the 19th and early 20th centuries, naturalists often collected bird's eggs and their nests. The practice of egg-collecting or oology is now illegal in many jurisdictions worldwide; the study of bird nests is called caliology.
Artificial bird nests
Bird nests are also built by humans to help in the conservation of certain birds. For example, artificial swallow nests are generally built with plaster, wood, terracotta or stucco.
Artificial nests, such as nest boxes, are an important conservation tool for many species, however nest box programs rarely compare their effectiveness with individuals not using nest boxes. Red-footed falcons using nest boxes in heavily managed landscapes produced fewer fledglings than those nesting in natural nests, but also than pairs nesting in nest boxes in more natural habitats.
| Biology and health sciences | Shelters and structures | Animals |
947692 | https://en.wikipedia.org/wiki/Dredging | Dredging | Dredging is the excavation of material from a water environment. Possible reasons for dredging include improving existing water features; reshaping land and water features to alter drainage, navigability, and commercial use; constructing dams, dikes, and other controls for streams and shorelines; and recovering valuable mineral deposits or marine life having commercial value. In all but a few situations the excavation is undertaken by a specialist floating plant, known as a dredger.
Usually the main objectives of dredging is to recover material of value, or to create a greater depth of water. Dredging systems can either be shore-based, brought to a location based on barges, or built into purpose-built vessels.
Dredging can have environmental impacts: it canc disturb marine sediments, creating dredge plumes which can lead to both short- and long-term water pollution, damage or destroy seabed ecosystems, and release legacy human-sourced toxins captured in the sediment. These environmental impacts can reduce marine wildlife populations, contaminate sources of drinking water, and interrupt economic activities such as fishing.
Description
Dredging is excavation carried out underwater or partially underwater, in shallow waters or ocean waters. It keeps waterways and ports navigable, and assists coastal protection, land reclamation and coastal redevelopment, by gathering up bottom sediments and transporting it elsewhere. Dredging can be done to recover materials of commercial value; these may be high value minerals or sediments such as sand and gravel that are used by the construction industry.
Dredging is a four-part process: loosening the material, bringing the material to the surface (together extraction), transportation and disposal.
The extract can be disposed of locally or transported by barge or in a liquid suspension in pipelines. Disposal can be to infill sites, or the material can be used constructively to replenish eroded sand that has been lost to coastal erosion, or constructively create sea-walls, building land or whole new landforms such as viable islands in coral atolls.
History
Ancient authors refer to harbour dredging. The seven arms of the Nile were channelled and wharfs built at the time of the pyramids (4000 BC), there was extensive harbour building in the eastern Mediterranean from 1000 BC and the disturbed sediment layers gives evidence of dredging. At Marseille, dredging phases are recorded from the third century BC onwards, the most extensive during the first century AD. The remains of three dredging boats have been unearthed; they were abandoned at the bottom of the harbour during the first and second centuries AD.
The Banu Musa brothers during the Muslim Golden Age in while working at the Bayt-Al-Hikmah (house of wisdom) in Baghdad, designed an original invention in their book named Book of Ingenious Devices, a grab machine that does not appear in any earlier Greek works. The grab they described was used to extract objects from underwater, and recover objects from the beds of streams.
During the renaissance Leonardo da Vinci drew a design for a drag dredger.
Dredging machines have been used during the construction of the Suez Canal from the late 1800s to present day expansions and maintenance. The completion of the Panama Canal in 1914, the most expensive U.S. engineering project at the time, relied extensively on dredging.
Purposes
Capital dredging: dredging carried out to create a new harbour, berth or waterway, or to deepen existing facilities in order to allow larger ships access. Because capital works usually involve hard material or high-volume works, the work is usually done using a cutter suction dredge or large trailing suction hopper dredge; but for rock works, drilling and blasting along with mechanical excavation may be used.
Land reclamation: dredging to mine sand, clay or rock from the seabed and using it to construct new land elsewhere. This is typically performed by a cutter-suction dredge or trailing suction hopper dredge. The material may also be used for flood or erosion control.
Maintenance: dredging to deepen or maintain navigable waterways or channels which are threatened to become silted with the passage of time, due to sedimented sand and mud, possibly making them too shallow for navigation. This is often carried out with a trailing suction hopper dredge. Most dredging is for this purpose, and it may also be done to maintain the holding capacity of reservoirs or lakes.
Harvesting materials: dredging sediment for elements like gold, diamonds or other valuable trace substances. Hobbyists examine their dredged matter to pick out items of potential value, similar to the hobby of metal detecting.
Fishing dredging is a technique for catching certain species of edible clams and crabs. In Louisiana and other American states, with salt water estuaries that can sustain bottom oyster beds, oysters are raised and harvested. A heavy rectangular metal scoop is towed astern of a moving boat with a chain bridle attached to a cable. This drags along the bottom scooping up oysters. It is periodically winched aboard and the catch is sorted and bagged for shipment.
Preparatory: dredging work and excavation for future bridges, piers or docks or wharves, This is often to build the foundations.
Winning construction materials: dredging sand and gravels from offshore licensed areas for use in construction industry, principally for use in concrete. This very specialist industry is focused in NW Europe, it uses specialized trailing suction hopper dredgers self discharging the dry cargo ashore. Land based old river beddings can be processed in this manner too.
Contaminant remediation: to reclaim areas affected by chemical spills, storm water surges (with urban runoff), and other soil contaminations, including silt from sewage sludge and from decayed matter, like wilted plants. Disposal becomes a proportionally large factor in these operations.
Flood prevention: dredging increases the channel depth and therefore increase a channel's capacity for carrying water.
Other
Beach nourishment: this is mining sand offshore and placing on a beach to replace sand eroded by storms or wave action. This enhances the recreational and protective function of the beach, which are also eroded by human activity. This is typically performed by a cutter-suction dredge or trailing suction hopper dredge.
Peat extraction: dredging poles or dredge hauls were used on the back of small boats to manually dredge the beds of peat-moor waterways. The extracted peat was used as a fuel. This tradition is now more or less obsolete. The tools are now significantly changed.
Removing rubbish and debris: often done in combination with maintenance dredging, this process removes non-natural matter from the bottoms of rivers and canals and harbours. Law enforcement agencies sometimes need to use a 'drag' to recover evidence or corpses from beneath the water.
Anti-eutrophication: A kind of contaminant remediation, dredging is an expensive option for the remediation of eutrophied (or de-oxygenated) water bodies; one of the causes is like mentioned above, sewage sludge. However, as artificially elevated phosphorus levels in the sediment aggravate the eutrophication process, controlled sediment removal is occasionally the only option for the reclamation of still waters.
Seabed mining: is a possible future use, recovering natural metal ore nodules from the sea's deepest troughs.
Types
Suction dredgers
These operate by sucking through a long tube like some vacuum cleaners but on a larger scale.
A plain suction dredger has no tool at the end of the suction pipe to disturb the material.
Trailing suction
A trailing suction hopper dredger (TSHD) trails its suction pipe when working. The pipe, which is fitted with a dredge drag head, loads the dredge spoil into one or more hoppers in the vessel.
When the hoppers are full, the TSHD sails to a disposal area and either dumps the material through doors in the hull or pumps the material out of the hoppers. Some dredges also self-offload using drag buckets and conveyors.
the largest trailing suction hopper dredgers in the world were Jan De Nul's Cristobal Colon (launched 4 July 2008) and her sister ship Leiv Eriksson (launched 4 September 2009). Main design specifications for the Cristobal Colon and the Leiv Eriksson are: 46,000 cubic metre hopper and a design dredging depth of 155 m. Next largest is HAM 318 (Van Oord) with its 37,293 cubic metre hopper and a maximum dredging depth of 101 m.
Cutter-suction
A cutter-suction dredger's (CSD) suction tube has a cutting mechanism at the suction inlet. The cutting mechanism loosens the bed material and transports it to the suction mouth. The dredged material is usually sucked up by a wear-resistant centrifugal pump and discharged either through a pipe line or to a barge. Cutter-suction dredgers are most often used in geological areas consisting of hard surface materials (for example gravel deposits or surface bedrock) where a standard suction dredger would be ineffective. They can, if sufficiently powerful, be used instead of underwater blasting.
, the most powerful cutter-suction dredger in the world is DEME's Spartacus, which entered service in 2021.
Auger suction
The auger dredge system functions like a cutter suction dredger, but the cutting tool is a rotating Archimedean screw set at right angles to the suction pipe. Mud Cat invented the auger dredge in the 1970s.
Jet-lift
These use the Venturi effect of a concentrated high-speed stream of water to pull the nearby water, together with bed material, into a pipe.
Air-lift
An airlift is a type of small suction dredge. It is sometimes used like other dredges. At other times, an airlift is handheld underwater by a diver. It works by blowing air into the pipe, and that air, being lighter than water, rises inside the pipe, dragging water with it.
Mechanical dredgers
Some bucket dredgers and grab dredgers are powerful enough to rip out coral to make a shipping channel through coral reefs.
Bucket dredgers
A bucket dredger is equipped with a bucket dredge, which is a device that picks up sediment by mechanical means, often with many circulating buckets attached to a wheel or chain.
Grab dredgers
A grab dredger picks up seabed material with a clam shell bucket, which hangs from an onboard crane or a crane barge, or is carried by a hydraulic arm, or is mounted like on a dragline. This technique is often used in excavation of bay mud. Most of these dredges are crane barges with spuds, steel piles that can be lowered and raised to position the dredge.
Backhoe/dipper dredgers
A backhoe/dipper dredger has a backhoe like on some excavators. A crude but usable backhoe dredger can be made by mounting a land-type backhoe excavator on a pontoon. The six largest backhoe dredgers in the world are currently the Vitruvius, the Mimar Sinan, Postnik Yakovlev (Jan De Nul), the Samson (DEME), the Simson and the Goliath (Van Oord). They featured barge-mounted excavators. Small backhoe dredgers can be track-mounted and work from the bank of ditches. A backhoe dredger is equipped with a half-open shell. The shell is filled moving towards the machine. Usually dredged material is loaded in barges. This machine is mainly used in harbours and other shallow water.
Excavator dredge attachments
The excavator dredge attachment uses the characteristics of cutter-suction dredgers, consisting of cutter heads and a suction pump for transferring material. These hydraulic attachments mount onto the boom arm of an excavator allowing an operator to maneuver the attachment along the shoreline and in shallow water for dredging.
Bed leveler
This is a bar or blade which is pulled over the seabed behind any suitable ship or boat. It has an effect similar to that of a bulldozer on land. The chain-operated steam dredger Bertha, built in 1844 to a design by Brunel and was the oldest operational steam vessel in Britain, was of this type.
Krabbelaar
This is an early type of dredger which was formerly used in shallow water in the Netherlands. It was a flat-bottomed boat with spikes sticking out of its bottom. As tide current pulled the boat, the spikes scraped seabed material loose, and the tide current washed the material away, hopefully to deeper water. Krabbelaar is the Dutch word for "scratcher".
Water injection
A water injection dredger uses a small jet to inject water under low pressure (to prevent the sediment from exploding into the surrounding waters) into the seabed to bring the sediment in suspension, which then becomes a turbidity current, which flows away down slope, is moved by a second burst of water from the WID or is carried away in natural currents. Water injection results in a lot of sediment in the water which makes measurement with most hydrographic equipment (for instance: singlebeam echosounders) difficult.
Pneumatic
These dredgers use a chamber with inlets, out of which the water is pumped with the inlets closed. It is usually suspended from a crane on land or from a small pontoon or barge. Its effectiveness depends on depth pressure.
Snagboat
A snagboat is designed to remove big debris such as dead trees and parts of trees from North America waterways.
Amphibious
Some of these are any of the above types of dredger, which can operate normally, or by extending legs, also known as spuds, so it stands on the seabed with its hull out of the water. Some forms can go on land.
Some of these are land-type backhoe excavators whose wheels are on long hinged legs so it can drive into shallow water and keep its cab out of water. Some of these may not have a floatable hull and, if so, cannot work in deep water. Oliver Evans (1755–1819) in 1804 invented the Oruktor Amphibolos, an amphibious dredger which was America's first steam-powered road vehicle.
Submersible
These are usually used to recover useful materials from the seabed. Many of them travel on continuous track. A unique variant is intended to walk on legs on the seabed.
Fishing
Fishing dredges are used to collect various species of clams, scallops, oysters or mussels from the seabed. Some dredges are also designed to catch crabs, sea urchins, sea cucumbers, and conch. These dredges have the form of a scoop made of chain mesh, and are towed by a fishing boat. Clam-specific dredges can utilize hydraulic injection to target deeper into the sand. Dredging can be destructive to the seabed and some scallop dredging has been replaced by collecting via scuba diving.
Notable individual dredgers
As of June 2018, the largest dredger in Asia is , a long dredger constructed in China, with a capacity of . An even larger dredger, retired in 1980, was the U.S. Army Corps of Engineers , which was long. The , a clamshell dredger that maintains levees in San Francisco Bay, has operated continuously since being built in 1936.
Dredge monitoring software
Dredgers are often equipped with dredge monitoring software to help the dredge operator position the dredger and monitor the current dredge level. The monitoring software often uses Real Time Kinematic satellite navigation to accurately record where the machine has been operating and to what depth the machine has dredged to.
Transportation and disposal of materials
In a "hopper dredger", the dredged materials end up in a large onboard hold called a "hopper." A suction hopper dredger is usually used for maintenance dredging. A hopper dredge usually has doors in its bottom to empty the dredged materials, but some dredges empty their hoppers by splitting the two-halves of their hulls on large hydraulic hinges. Either way, as the vessel dredges, excess water in the dredged materials is spilled off as the heavier solids settle to the bottom of the hopper. This excess water is returned to the sea to reduce weight and increase the amount of solid material (or slurry) that can be carried in one load. When the hopper is filled with slurry, the dredger stops dredging and goes to a dump site and empties its hopper.
Some hopper dredges are designed so they can also be emptied from above using pumps if dump sites are unavailable or if the dredge material is contaminated. Sometimes the slurry of dredgings and water is pumped straight into pipes which deposit it on nearby land. These pipes are also commonly known as dredge hoses, too. There are a few different types of dredge hoses that differ in terms of working pressure, float-ability, armored or not etc. Suction hoses, discharge armored hoses and self-floating hoses are some of the popular types engineered for transporting and discharging dredge materials. Some even had the pipes or hoses customised to exact dredging needs etc. Other times, it is pumped into barges (also called scows), which deposit it elsewhere while the dredge continues its work.
A number of vessels, notably in the UK and NW Europe de-water the hopper to dry the cargo to enable it to be discharged onto a quayside 'dry'. This is achieved principally using self discharge bucket wheel, drag scraper or excavator via conveyor systems.
When contaminated (toxic) sediments are to be removed, or large volume inland disposal sites are unavailable, dredge slurries are reduced to dry solids via a process known as dewatering, to minimize the dredge plume. Current dewatering techniques employ either centrifuges, geotube containers, large textile based filters or polymer flocculant/congealant based apparatus.
In many projects, slurry dewatering is performed in large inland settling pits, although this is becoming less and less common as mechanical dewatering techniques continue to improve.
Similarly, many groups (most notable in east Asia) are performing research towards utilizing dewatered sediments for the production of concretes and construction block, although the high organic content (in many cases) of this material is a hindrance toward such ends.
The proper management of contaminated sediments is a modern-day issue of significant concern. Because of a variety of maintenance activities, thousands of tonnes of contaminated sediment are dredged worldwide from commercial ports and other aquatic areas at high level of industrialization. Dredged material can be reused after appropriate decontamination. A variety of processes has been proposed and tested at different scales of application (technologies for environmental remediation). Once decontaminated, the material could well suit the building industry, or could be used for beach nourishment.
Environmental impacts
Dredging can disturb aquatic ecosystems, often with adverse impacts. In addition, dredge spoils may contain toxic chemicals that may have an adverse effect on the disposal area; furthermore, the process of dredging often dislodges chemicals residing in benthic substrates and injects them into the water column, where they become toxic dredge plumes.
Dredging can have numerous significant impacts on the environment, including the following:
Release of toxic chemicals (including heavy metals and PCB) from bottom sediments into the water column.
Short term increases in turbidity, which can affect aquatic species metabolism and interfere with spawning. Suction dredging activity is allowed only during non-spawning time frames set by fish and game (in-water work periods).
Secondary impacts to marsh productivity from sedimentation and general changes in wetland chemistry after dredging.
Tertiary impacts to avifauna which may prey upon contaminated aquatic organisms.
Secondary impacts to aquatic and benthic organisms' metabolism and mortality.
Possible contamination of dredge spoils sites.
Changes to the topography by creating "spoil islands" from the accumulated spoil.
Releases toxic compound Tributyltin, a biocide often used in anti-fouling paint banned in 2008, into the water.
The nature of dredging operations and possible environmental impacts requires that the activity often be closely regulated and requires comprehensive regional environmental impact assessments alongside continuous monitoring. For example, in the U.S., the Clean Water Act requires that any discharge of dredged or fill materials into "waters of the United States," including wetlands, is forbidden unless authorized by a permit issued by the Army Corps of Engineers. Due to potential environmental impacts, dredging is often restricted to licensed areas, with vessel activity monitored closely using automatic GPS systems.
Major dredging companies
According to a Rabobank outlook report in 2013, the largest dredging companies in the world are in order of size, based on dredging sales in 2012
China Harbour Engineering (China)
Jan De Nul (Belgium)
DEME (Belgium)
Royal Boskalis Westminster (Netherlands)
Van Oord Dredging and Marine Contractors (Netherlands)
National Marine Dredging Company (United Arab Emirates)
Great Lakes Dredge & Dock Company (United States)
Notable dredging companies in North America
Manson Construction Co. (United States)
Notable dredging companies in South Asia
Dredging Corporation of India
Adani Ports & SEZ (India)
Maldives Transport and Contracting Company (Maldives)
Images
| Technology | Hydraulic infrastructure | null |
948422 | https://en.wikipedia.org/wiki/Fell%20Terrier | Fell Terrier | Fell terrier refers to a regional type of long-legged working terrier, not a specific breed of dog.
Description and purpose
Fell terriers are types of small working terriers developed in the fell country of Northern England and used as hunting dogs. They may be crossbred or purebred. Fell terrier types are typically small, usually 10- 15 lbs/6.5 kg, and with a narrow chest, so as to fit into the tunnels of the animals they hunt. Fell terriers are long-legged, with a rough textured coat, often red or black in colour. The tail traditionally is docked; in the United States the tail is not required to be docked. Crossbreeding with other hunting terriers in the beginning caused the appearance to vary.
Fell terriers are bred for hunting ability and gameness rather than to a standard of appearance (breed type). They hunt in packs or alone. The fell terrier was originally developed by Ullswater Hunt Master Joe Bowman, an early Border Terrier breeder, where he used the best red fell terriers available to him, so that he could continue his efforts to refine the fell terrier even further to hunt the large fell fox that was believed to cause serious losses for sheep stockmen. The dog needed long legs to follow hunters through heavy snow, and a narrow chest to follow the fox in a stony underground den. In the hunt, a terrier follows the red fox underground into its den, where it either kills the fox, bolts it or holds it until the hunter (terrierman) digs the dog and fox up. The original fell terrier bloodlines extend down from Harry Hardisty Turk and Sid Wilkinson's Rock, Wilkinson's Rock the most important stud dog of his Era bred down from Fred Barker and Anthony Barker's chowt-face rock terriers. Others who have contributed are Garry Middleton,
Brian Nuttall, John Park, Ken Gould, Frank Buck, Cyril Breay, Joe Armstrong, Anthony Barker, Fred Barker, Maurice Bell, Anthony Chapman, John Cowen, Tommy Dobson, Eddie Pool, Graham Ward, and Sid Wilkinson. Fell terriers have been used in the United States for several generations hunting small game and have been known under the name Patterdale terrier. For the most part remain unchanged by the hunter who keeps the standard.
Breeds
Several named breeds have been developed from the fell terrier type, such as the Border Terrier, Lakeland Terrier, Patterdale Terrier, scorrier terrier, Welsh Terrier, and other locally developed breeds. All are sometimes called "fell terrier" interchangeably with their breed name. The "National Terriers Club LLC", The "American Fell Terrier International" has published a Fell Terrier standard. In Germany, the Jagdterrier was developed out of Fell terrier stock in the 1920s. Fell terriers may be descended from a very old type of long-legged terrier referred to as the rough-coated Black and Tan, similar to today's Welsh Terrier.
| Biology and health sciences | Dogs | Animals |
948821 | https://en.wikipedia.org/wiki/Afrocarpus | Afrocarpus | Afrocarpus is a genus of conifer of the family Podocarpaceae. Two to six species are recognized. They are evergreen trees native to Africa. Afrocarpus was designated a genus in 1989, when several species formerly classified in Podocarpus and Nageia were reclassified.
Taxonomy
Afrocarpus gaussenii was based on a single specimen of a cultivated individual of Afrocarpus falcatus in Madagascar. Its distinctive features might have resulted from the conditions of its cultivation. No species of Afrocarpus is known to be native to Madagascar.
In a recent treatment of Afrocarpus, only two species were recognized; A. dawei, A. gracilior, and A. usambarensis were sunk into A. falcatus. The reason for this merger was that "variation across the group appears to be essentially continuous".
Studies based on anatomical, biogeographical, morphological, and DNA evidence suggest the following relationships:
Species
Description
Afrocarpus are evergreen trees. The individuals of the largest species, Afrocarpus falcatus, may reach a height of 60 meters. The thin bark often peels with scale-like plates.
The leaves are simple and flat. The phyllotaxis or leaf arrangement is usually spiral but may be opposite on young plants. The leaves are generally lanceolate in shape and coriaceous in texture. They have a single visible midrib. Stomata are found on both surfaces of the leaf.
Afrocarpus are dioecious, with male pollen cones and female seed cones borne on separate individual plants. The cones are short pedunculate and usually develop from axillary buds.
The male pollen cones are narrowly cylindrical and resemble catkins. They grow in small groups of two or three cones. The peduncles are glabrous. Each pollen cone has numerous spirally inserted microsporophylls each with two basal pollen sacs producing bisaccate pollen.
The female seed cones are solitary. Their peduncles may have small scale leaves. The cones consist of several sterile cone scales and one fertile cone scale with just one seed producing ovule. The sterile scales wither as the cone matures, unlike in the closely related genus Podocarpus where the scales fuse to form a fleshy receptacle. A part of the scale supporting the ovule develops into a rounded fleshy covering enclosing the seed entirely known as the epimatium. At maturity the epimatium varies in shape from subglobose to elliptic or obovoid and in color from greenish to yellow or brown.
Distribution
As the name intimates, Afrocarpus is native to Africa. The species are distributed through the Afromontane forests of eastern and southern Africa, descending to the Indian Ocean coast in South Africa. The genus is native to Burundi, Democratic Republic of the Congo, Eswatini, Ethiopia, Kenya, Malawi, Mozambique, Rwanda, São Tomé and Príncipe, South Africa, Tanzania and Uganda.
The podocarps are associated with the ancient supercontinent of Gondwana, where they were characteristic of the cool, moist southern Gondwana flora. Gondwana broke up into the continents of South America, Africa, India, Australia, and Antarctica between 160 and 30 million years ago. As Africa drifted north, it became hotter and drier, and the podocarps generally retreated to the cool, moist highlands of eastern and southern Africa.
Uses
In South Africa, this wood is mostly used to make exclusive furniture.
| Biology and health sciences | Pinophyta (Conifers) | Plants |
948919 | https://en.wikipedia.org/wiki/Goliath%20birdeater | Goliath birdeater | The Goliath birdeater (Theraphosa blondi) belongs to the tarantula family Theraphosidae. Found in northern South America, it is the largest spider in the world by mass () and body length (up to ), and second to the giant huntsman spider by leg span. It is also called the Goliath tarantula or Goliath bird-eating spider; the practice of calling theraphosids "bird-eating" derives from an early 18th-century copper engraving by Maria Sibylla Merian that shows one eating a hummingbird. Despite the spider's name, it rarely preys on birds.
Characteristics
These spiders can have a leg span of up to , a body length of up to , and can weigh up to . Birdeaters are one of the few tarantula species that lack tibial spurs, located on the first pair of legs of most adult males. They are mostly tan to light brown and golden-hued.
Life cycle
Unlike other species of spider/tarantula, females rarely eat the males during mating. Females mature in 3–6 years and have an average lifespan of 15 to 25 years. Males die soon after maturity and have a lifespan of three to six years. Colors range from dark to light brown with faint markings on the legs. Bird-eaters have hair on their bodies, abdomens, and legs. The female lays 100 to 200 eggs, which hatch into spiderlings within 6–8 weeks.
Behaviour
Defenses
In response to threats, Goliath birdeaters stridulate by rubbing setae on their pedipalps and legs. Also, when threatened they rub their abdomen with their hind legs and release hairs that are a severe irritant to the skin and mucous membranes. These urticating hairs can be harmful to humans.
Like all tarantulas, T. blondi spiders have fangs large enough () to break the skin of a human. They carry venom in their fangs and have been known to bite when threatened, but the venom is relatively harmless and its effects are comparable to those of a wasp's sting. Tarantulas generally bite humans only in self-defense, and these bites do not always result in envenomation (known as a "dry bite”).
Feeding
Despite its name, the Goliath birdeater only rarely actually preys on birds; in the wild, its diet consists primarily of other large arthropods, worms, and amphibians. However, because of its size and opportunistic predatory behavior, this species commonly kills and consumes a variety of insects and small terrestrial vertebrates. They do not consume their prey in the open; rather, they drag it back to their burrow and begin the digesting process. They do this by liquefying the insides of their prey and proceed to suck it dry. In the wild, T. blondi has been observed feeding on rodents, frogs, toads, lizards, and even snakes.
Distribution and habitat
The Goliath birdeater is native to the upland rainforest regions of Northern South America: Suriname, Guyana, French Guiana, northern Brazil, eastern Colombia, and southern Venezuela. Most noticeable in the Amazon rainforest, the spider is terrestrial, living in deep burrows, and is found commonly in marshy or swampy areas. It is a nocturnal species.
Culinary use
The Goliath birdeater is an edible spider. The spider is part of the local cuisine in northeastern South America, prepared by singeing off the urticating hairs and roasting it in banana leaves. The flavor has been described as "shrimp-like".
| Biology and health sciences | Spiders | Animals |
2095010 | https://en.wikipedia.org/wiki/Merbromin | Merbromin | Merbromin (marketed as Mercurochrome, Merbromine, Mercurocol, Sodium mercurescein, Asceptichrome, Supercrome, Brocasept and Cinfacromin) is an organomercuric disodium salt compound used as a topical antiseptic for minor cuts and scrapes and as a biological dye. Readily available in most countries, it is no longer sold in Switzerland, Brazil, France, Iran, Germany, Denmark, or the United States, due to its mercury content.
Uses
Merbromin's best-known use is as a topical antiseptic to treat minor wounds, burns, and scratches. It is also used in the antisepsis of the umbilical cord, and the antisepsis of wounds with inhibited scar formation, such as neuropathic ulcers and diabetic foot sores. When applied on a wound, it stains the skin a distinctive carmine red, which can persist through repeated washings. Due to its persistence and to its lethality to bacteria, Merbromin is useful on infections of the fingernail or toenail.
In 1998, the U.S. Food and Drug Administration reclassified merbromin from "generally recognized as safe" to "untested," due to a lack of recent studies or updated supporting information. Consequently, its use in the United States has been superseded by other agents (e.g., povidone iodine, benzalkonium chloride, chloroxylenol).
Synthesis
Merbromin is synthesized by combining dibromofluorescein with mercuric acetate and sodium hydroxide or, alternatively, through action of the mercuric acetate upon (or combining with) sodium dibromofluorescein. Because of its anionic character, it is chemically incompatible with acids, the majority of alkaloid salts and most local anesthetics.
Mercurochrome
Merbromin is sold under the trade name Mercurochrome (in which the suffix "-chrome" denotes "color"). The name is also commonly used for over-the-counter antiseptic solutions consisting of merbromin (typically at 2% concentration) dissolved in either ethyl alcohol (tincture) or water (aqueous).
Its antiseptic qualities were discovered in 1918 by Hugh H. Young, a physician at Johns Hopkins Hospital. The chemical soon became popular among parents and physicians for everyday antiseptic uses, in part because the dye component made it easy to see where the antiseptic had been applied.
On 19 October 1998, citing potential for mercury poisoning, the US Food and Drug Administration (FDA) reclassified merbromin from "generally recognized as safe" to "untested," effectively halting its distribution within the United States. Sales were subsequently halted in Brazil (2001), Germany (2003), and France (2006). It remains readily available in most other countries.
Within the United States, products such as Humco Mercuroclear ("Aqueous solution of benzalkonium chloride and lidocaine hydrochloride") play on the brand recognition history of Mercurochrome but substitute other ingredients with similar properties. In Canada, Jean Coutu Group markets a chlorhexidine solution under the name Mercurochrome.
| Physical sciences | Organic salts | Chemistry |
2096232 | https://en.wikipedia.org/wiki/Geology%20of%20the%20Moon | Geology of the Moon | The geology of the Moon (sometimes called selenology, although the latter term can refer more generally to "lunar science") is quite different from that of Earth. The Moon lacks a true atmosphere, and the absence of free oxygen and water eliminates erosion due to weather. Instead, the surface is eroded much more slowly through the bombardment of the lunar surface by micrometeorites. It does not have any known form of plate tectonics, it has a lower gravity, and because of its small size, it cooled faster. In addition to impacts, the geomorphology of the lunar surface has been shaped by volcanism, which is now thought to have ended less than 50 million years ago. The Moon is a differentiated body, with a crust, mantle, and core.
Geological studies of the Moon are based on a combination of Earth-based telescope observations, measurements from orbiting spacecraft, lunar samples, and geophysical data. Six locations were sampled directly during the crewed Apollo program landings from 1969 to 1972, which returned of lunar rock and lunar soil to Earth In addition, three robotic Soviet Luna spacecraft returned another of samples, and the Chinese robotic Chang'e 5 returned a sample of 1,731 g (61.1 oz) in 2020.
The Moon is the only extraterrestrial body for which we have samples with a known geologic context. A handful of lunar meteorites have been recognized on Earth, though their source craters on the Moon are unknown. A substantial portion of the lunar surface has not been explored, and a number of geological questions remain unanswered.
Elemental composition
Elements known to be present on the lunar surface include, among others, oxygen (O), silicon (Si), iron (Fe), magnesium (Mg), calcium (Ca), aluminium (Al), manganese (Mn) and titanium (Ti). Among the more abundant are oxygen, iron and silicon. The oxygen content is estimated at 45% (by weight). Carbon (C) and nitrogen (N) appear to be present only in trace quantities from deposition by solar wind.
Formation
For a long period of time, the fundamental question regarding the history of the Moon was of its origin. Early hypotheses included fission from Earth, capture, and co-accretion. Today, the giant-impact hypothesis is widely accepted by the scientific community.
Geologic history
The geological history of the Moon has been defined into six major epochs, called the lunar geologic timescale. Starting about 4.5 billion years ago, the newly formed Moon was in a molten state and was orbiting much closer to Earth resulting in tidal forces. These tidal forces deformed the molten body into an ellipsoid, with the major axis pointed towards Earth.
The first important event in the geologic evolution of the Moon was the crystallization of the near global magma ocean. It is not known with certainty what its depth was, but several studies imply a depth of about 500 km or greater. The first minerals to form in this ocean were the iron and magnesium silicates olivine and pyroxene. Because these minerals were denser than the molten material around them, they sank. After crystallization was about 75% complete, less dense anorthositic plagioclase feldspar crystallized and floated, forming an anorthositic crust about 50 km in thickness. The majority of the magma ocean crystallized quickly (within about 100 million years or less), though the final remaining KREEP-rich magmas, which are highly enriched in incompatible and heat-producing elements, could have remained partially molten for several hundred million (or perhaps 1 billion) years. It appears that the final KREEP-rich magmas of the magma ocean eventually became concentrated within the region of Oceanus Procellarum and the Imbrium basin, a unique geologic province that is now known as the Procellarum KREEP Terrane.
Quickly after the lunar crust formed, or even as it was forming, different types of magmas that would give rise to the Mg-suite norites and troctolites began to form, although the exact depths at which this occurred are not known precisely. Recent theories suggest that Mg-suite plutonism was largely confined to the region of the Procellarum KREEP Terrane, and that these magmas are genetically related to KREEP in some manner, though their origin is still highly debated in the scientific community. The oldest of the Mg-suite rocks have crystallization ages of about 3.85 Ga. However, the last large impact that could have been excavated deep into the crust (the Imbrium basin) also occurred at 3.85 Ga before present. Thus, it seems probable that Mg-suite plutonic activity continued for a much longer time, and that younger plutonic rocks exist deep below the surface.
Analysis of the samples from the Moon seems to show that a lot of the Moon's impact basins formed in a short amount of time between about 4 and 3.85 Ga ago. This hypothesis is referred to as the lunar cataclysm or late heavy bombardment. However, it is now recognized that ejecta from the Imbrium impact basin (one of the youngest large impact basins on the Moon) should be found at all of the Apollo landing sites. It is thus possible that ages for some impact basins (in particular Mare Nectaris) could have been mistakenly assigned the same age as Imbrium.
The lunar maria represent ancient flood basaltic eruptions. In comparison to terrestrial lavas, these contain higher iron abundances, have low viscosities, and some contain highly elevated abundances of the titanium-rich mineral ilmenite. The majority of basaltic eruptions occurred between about 3 and 3.5 Ga ago, though some mare samples have ages as old as 4.2 Ga. The youngest (based on the method of crater counting) was long thought to date to 1 billion years ago, but research in the 2010s has found evidence of eruptions from less than 50 million years in the past. Along with mare volcanism came pyroclastic eruptions, which launched molten basaltic materials hundreds of kilometers away from the volcano. A large portion of the mare formed, or flowed into, the low elevations associated with the nearside impact basins. However, Oceanus Procellarum does not correspond to any known impact structure, and the lowest elevations of the Moon within the farside South Pole-Aitken basin are only modestly covered by mare (see lunar mare for a more detailed discussion).
Impacts by meteorites and comets are the only abrupt geologic force acting on the Moon today, though the variation of Earth tides on the scale of the Lunar anomalistic month causes small variations in stresses. Some of the most important craters used in lunar stratigraphy formed in this recent epoch. For example, the crater Copernicus, which has a depth of 3.76 km and a radius of 93 km, is estimated to have formed about 900 million years ago (though this is debatable). The Apollo 17 mission landed in an area in which the material coming from the crater Tycho might have been sampled. The study of these rocks seem to indicate that this crater could have formed 100 million years ago, though this is debatable as well. The surface has also experienced space weathering due to high energy particles, solar wind implantation, and micrometeorite impacts. This process causes the ray systems associated with young craters to darken until it matches the albedo of the surrounding surface. However, if the composition of the ray is different from the underlying crustal materials (as might occur when a "highland" ray is emplaced on the mare), the ray could be visible for much longer times.
After resumption of Lunar exploration in the 1990s, it was discovered there are scarps across the globe that are caused by the contraction due to cooling of the Moon.
Strata and epochs
At the top of the Moon’s stratigraphy is the Copernican unit consisting of craters with a ray system. Below this is the Eratosthenian unit, defined by craters with established impact crater morphology, but lacking the ray system of the Copernican. These two units are present in smaller spots on the lunar surface. Further down the stratigraphy are the Mare units (previously known as the Procellarian unit), and the Imbrian unit which is related to ejecta and tectonics from the Imbrium basin. The bottom of the lunar stratigraphy is the pre-Nectarian unit, which consists of old crater plains.
Lunar landscape
The lunar landscape is characterized by impact craters, their ejecta, a few volcanoes, hills, lava flows and depressions filled by lava.
Highlands
The most distinctive aspect of the Moon is the contrast between its bright and dark zones. Lighter surfaces are the lunar highlands, which receive the name of terrae (singular terra, from the Latin for earth, land), and the darker plains are called maria (singular mare, from the Latin for sea), after Johannes Kepler who introduced the names in the 17th century. The highlands are anorthositic in composition, whereas the maria are basaltic. The maria often coincide with the "lowlands," but the lowlands (such as within the South Pole-Aitken basin) are not always covered by maria. The highlands are older than the visible maria, and hence are more heavily cratered.
Maria
The major products of volcanic processes on the Moon are evident to Earth-bound observers in the form of the lunar maria. These are large flows of basaltic lava that correspond to low-albedo surfaces covering nearly a third of the near side. Only a few percent of the farside has been affected by mare volcanism. Even before the Apollo missions confirmed it, most scientists already thought that the maria are lava-filled plains, because they have lava flow patterns and collapses attributed to lava tubes.
The ages of the mare basalts have been determined both by direct radiometric dating and by the technique of crater counting. The oldest radiometric ages are about 4.2 Ga (billion years), and ages of most of the youngest maria lavas have been determined from crater counting to be about 1 Ga. Due to better resolution of more recent imagery, about 70 small areas called irregular mare patches (each area only a few hundred meters or a few kilometers across) have been found in the maria that crater counting suggests were sites of volcanic activity in the geologically much more recent past (less than 50 million years). Volumetrically, most of the mare formed between about 3 and 3.5 Ga before present. The youngest lavas erupted within Oceanus Procellarum, whereas some of the oldest appear to be located on the farside. The maria are clearly younger than the surrounding highlands given their lower density of impact craters.
A large portion of maria erupted within, or flowed into, the low-lying impact basins on the lunar nearside. However, it is unlikely that a causal relationship exists between the impact event and mare volcanism because the impact basins are much older (by about 500 million years) than the mare fill. Furthermore, Oceanus Procellarum, which is the largest expanse of mare volcanism on the Moon, does not correspond to any known impact basin. It is commonly suggested that the reason the mare only erupted on the nearside is that the nearside crust is thinner than the farside. Although variations in the crustal thickness might act to modulate the amount of magma that ultimately reaches the surface, this hypothesis does not explain why the farside South Pole-Aitken basin, whose crust is thinner than Oceanus Procellarum, was only modestly filled by volcanic products.
Another type of deposit associated with the maria, although it also covers the highland areas, are the "dark mantle" deposits. These deposits cannot be seen with the naked eye, but they can be seen in images taken from telescopes or orbiting spacecraft. Before the Apollo missions, scientists predicted that they were deposits produced by pyroclastic eruptions. Some deposits appear to be associated with dark elongated ash cones, reinforcing the idea of pyroclasts. The existence of pyroclastic eruptions was later confirmed by the discovery of glass spherules similar to those found in pyroclastic eruptions here on Earth.
Many of the lunar basalts contain small holes called vesicles, which were formed by gas bubbles exsolving from the magma at the vacuum conditions encountered at the surface. It is not known with certainty which gases escaped these rocks, but carbon monoxide is one candidate.
The samples of pyroclastic glasses are of green, yellow, and red tints. The difference in color indicates the concentration of titanium that the rock has, with the green particles having the lowest concentrations (about 1%), and red particles having the highest concentrations (up to 14%, much more than the basalts with the highest concentrations).
Rilles
Rilles on the Moon sometimes resulted from the formation of localized lava channels. These generally fall into three categories, consisting of sinuous, arcuate, or linear shapes. By following these meandering rilles back to their source, they often lead to an old volcanic vent. One of the most notable sinuous rilles is the Vallis Schröteri feature, located in the Aristarchus plateau along the eastern edge of Oceanus Procellarum. An example of a sinuous rille exists at the Apollo 15 landing site, Rima Hadley, located on the rim of the Imbrium Basin. Based on observations from the mission, it is generally thought that this rille was formed by volcanic processes, a topic long debated before the mission took place.
Domes
A variety of shield volcanoes can be found in selected locations on the lunar surface, such as on Mons Rümker. These are thought to be formed by relatively viscous, possibly silica-rich lava, erupting from localized vents. The resulting lunar domes are wide, rounded, circular features with a gentle slope rising in elevation a few hundred meters to the midpoint. They are typically 8–12 km in diameter, but can be up to 20 km across. Some of the domes contain a small pit at their peak.
Wrinkle ridges
Wrinkle ridges are features created by compressive tectonic forces within the maria. These features represent buckling of the surface and form long ridges across parts of the maria. Some of these ridges may outline buried craters or other features beneath the maria. A prime example of such an outlined feature is the crater Letronne.
Grabens
Grabens are tectonic features that form under extensional stresses. Structurally, they are composed of two normal faults, with a down-dropped block between them. Most grabens are found within the lunar maria near the edges of large impact basins.
Impact craters
The origin of the Moon's craters as impact features became widely accepted only in the 1960s. This realization allowed the impact history of the Moon to be gradually worked out by means of the geologic principle of superposition. That is, if a crater (or its ejecta) overlaid another, it must be the younger. The amount of erosion experienced by a crater was another clue to its age, though this is more subjective. Adopting this approach in the late 1950s, Gene Shoemaker took the systematic study of the Moon away from the astronomers and placed it firmly in the hands of the lunar geologists.
Impact cratering is the most notable geological process on the Moon. The craters are formed when a solid body, such as an asteroid or comet, collides with the surface at a high velocity (mean impact velocities for the Moon are about 17 km per second). The kinetic energy of the impact creates a compression shock wave that radiates away from the point of entry. This is succeeded by a rarefaction wave, which is responsible for propelling most of the ejecta out of the crater. Finally there is a hydrodynamic rebound of the floor that can create a central peak.
These craters appear in a continuum of diameters across the surface of the Moon, ranging in size from tiny pits to the immense South Pole–Aitken basin with a diameter of nearly 2,500 km and a depth of 13 km. In a very general sense, the lunar history of impact cratering follows a trend of decreasing crater size with time. In particular, the largest impact basins were formed during the early periods, and these were successively overlaid by smaller craters. The size frequency distribution (SFD) of crater diameters on a given surface (that is, the number of craters as a function of diameter) approximately follows a power law with increasing number of craters with decreasing crater size. The vertical position of this curve can be used to estimate the age of the surface.
The most recent impacts are distinguished by well-defined features, including a sharp-edged rim. Small craters tend to form a bowl shape, whereas larger impacts can have a central peak with flat floors. Larger craters generally display slumping features along the inner walls that can form terraces and ledges. The largest impact basins, the multiring basins, can even have secondary concentric rings of raised material.
The impact process excavates high albedo materials that initially gives the crater, ejecta, and ray system a bright appearance. The process of space weathering gradually decreases the albedo of this material such that the rays fade with time. Gradually the crater and its ejecta undergo impact erosion from micrometeorites and smaller impacts. This erosional process softens and rounds the features of the crater. The crater can also be covered in ejecta from other impacts, which can submerge features and even bury the central peak.
The ejecta from large impacts can include large blocks of material that reimpact the surface to form secondary impact craters. These craters are sometimes formed in clearly discernible radial patterns, and generally have shallower depths than primary craters of the same size. In some cases an entire line of these blocks can impact to form a valley. These are distinguished from catena, or crater chains, which are linear strings of craters that are formed when the impact body breaks up prior to impact.
Generally speaking, a lunar crater is roughly circular in form. Laboratory experiments at NASA's Ames Research Center have demonstrated that even very low-angle impacts tend to produce circular craters, and that elliptical craters start forming at impact angles below five degrees. However, a low angle impact can produce a central peak that is offset from the midpoint of the crater. Additionally, the ejecta from oblique impacts show distinctive patterns at different impact angles: asymmetry starting around 60˚ and a wedge-shaped "zone of avoidance" free of ejecta in the direction the projectile came from starting around 45˚.
Dark-halo craters are formed when an impact excavates lower albedo material from beneath the surface, then deposits this darker ejecta around the main crater. This can occur when an area of darker basaltic material, such as that found on the maria, is later covered by lighter ejecta derived from more distant impacts in the highlands. This covering conceals the darker material below, which is later excavated by subsequent craters.
The largest impacts produced melt sheets of molten rock that covered portions of the surface that could be as thick as a kilometer. Examples of such impact melt can be seen in the northeastern part of the Mare Orientale impact basin.
Regolith
The surface of the Moon has been subject to billions of years of collisions with both small and large asteroidal and cometary materials. Over time, these impact processes have pulverized and "gardened" the surface materials, forming a fine-grained layer termed regolith. The thickness of the lunar regolith varies between beneath the younger maria, to up to beneath the oldest surfaces of the lunar highlands. The regolith is predominantly composed of materials found in the region, but also contains traces of materials ejected by distant impact craters. The term mega-regolith is often used to describe the heavily fractured bedrock directly beneath the near-surface regolith layer.
The regolith contains rocks, fragments of minerals from the original bedrock, and glassy particles formed during the impacts. In most of the lunar regolith, half of the particles are made of mineral fragments fused by the glassy particles; these objects are called agglutinates. The chemical composition of the regolith varies according to its location; the regolith in the highlands is rich in aluminium and silica, just as the rocks in those regions. The regolith in the maria is rich in iron and magnesium and is silica-poor, as are the basaltic rocks from which it is formed.
The lunar regolith is very important because it also stores information about the history of the Sun. The atoms that compose the solar wind – mostly hydrogen, helium, neon, carbon and nitrogen – hit the lunar surface and insert themselves into the mineral grains. Upon analyzing the composition of the regolith, particularly its isotopic composition, it is possible to determine if the activity of the Sun has changed with time. The gases of the solar wind could be useful for future lunar bases, because oxygen, hydrogen (water), carbon and nitrogen are not only essential to sustain life, but are also potentially very useful in the production of fuel. The composition of the lunar regolith can also be used to infer its source origin.
Lunar lava tubes
Lunar lava tubes form a potentially important location for constructing a future lunar base, which may be used for local exploration and development, or as a human outpost to serve exploration beyond the Moon. A lunar lava cave potential has long been suggested and discussed in literature and thesis. Any intact lava tube on the Moon could serve as a shelter from the severe environment of the lunar surface, with its frequent meteorite impacts, high-energy ultraviolet radiation and energetic particles, and extreme diurnal temperature variations. Following the launch of the Lunar Reconnaissance Orbiter, many lunar lava tubes have been imaged. These lunar pits are found in several locations across the Moon, including Marius Hills, Mare Ingenii and Mare Tranquillitatis.
Lunar magma ocean
The first rocks brought back by Apollo 11 were basalts. Although the mission landed on Mare Tranquillitatis, a few millimetric fragments of rocks coming from the highlands were picked up. These are composed mainly of plagioclase feldspar; some fragments were composed exclusively of anorthite. The identification of these mineral fragments led to the bold hypothesis that a large portion of the Moon was once molten, and that the crust formed by fractional crystallization of this magma ocean.
A natural outcome of the hypothetical giant-impact event is that the materials that re-accreted to form the Moon must have been hot. Current models predict that a large portion of the Moon would have been molten shortly after the Moon formed, with estimates for the depth of this magma ocean ranging from about 500 km to complete melting. Crystallization of this magma ocean would have given rise to a differentiated body with a compositionally distinct crust and mantle and accounts for the major suites of lunar rocks.
As crystallization of the lunar magma ocean proceeded, minerals such as olivine and pyroxene would have precipitated and sank to form the lunar mantle. After crystallization was about three-quarters complete, anorthositic plagioclase would have begun to crystallize, and because of its low density, float, forming an anorthositic crust. Importantly, elements that are incompatible (i.e., those that partition preferentially into the liquid phase) would have been progressively concentrated into the magma as crystallization progressed, forming a KREEP-rich magma that initially should have been sandwiched between the crust and mantle. Evidence for this scenario comes from the highly anorthositic composition of the lunar highland crust, as well as the existence of KREEP-rich materials. Additionally, zircon analysis of Apollo 14 samples suggests the lunar crust differentiated 4.51±0.01 billion years ago.
Lunar rocks
Surface materials
The Apollo program brought back of lunar surface material, most of which is stored at the Lunar Receiving Laboratory in Houston, Texas, and the uncrewed Soviet Luna programme returned of lunar material. These rocks have proved to be invaluable in deciphering the geologic evolution of the Moon. Lunar rocks are in large part made of the same common rock forming minerals as found on Earth, such as olivine, pyroxene, and plagioclase feldspar (anorthite). Plagioclase feldspar is mostly found in the lunar crust, whereas pyroxene and olivine are typically seen in the lunar mantle. The mineral ilmenite is highly abundant in some mare basalts, and a new mineral named armalcolite (named for Armstrong, Aldrin, and Collins, the three members of the Apollo 11 crew) was first discovered in the lunar samples.
The maria are composed predominantly of basalt, whereas the highland regions are iron-poor and composed primarily of anorthosite, a rock composed primarily of calcium-rich plagioclase feldspar. Another significant component of the crust are the igneous Mg-suite rocks, such as the troctolites, norites, and KREEP-basalts. These rocks are thought to be related to the petrogenesis of KREEP.
Composite rocks on the lunar surface often appear in the form of breccias. Of these, the subcategories are called fragmental, granulitic, and impact-melt breccias, depending on how they were formed. The mafic impact melt breccias, which are typified by the low-K Fra Mauro composition, have a higher proportion of iron and magnesium than typical upper crust anorthositic rocks, as well as higher abundances of KREEP.
Composition of the maria
The main characteristics of the basaltic rocks with respect to the rocks of the lunar highlands is that the basalts contain higher abundances of olivine and pyroxene, and less plagioclase. They are richer in iron than terrestrial basalts, and also have lower viscosities. Some of them have high abundances of a ferro-titanic oxide called ilmenite. Because the first sampling of rocks contained a high content of ilmenite and other related minerals, they received the name of "high titanium" basalts. The Apollo 12 mission returned to Earth with basalts of lower titanium concentrations, and these were dubbed "low titanium" basalts. Subsequent missions, including the Soviet robotic probes, returned with basalts with even lower concentrations, now called "very low titanium" basalts. The Clementine space probe returned data showing that the mare basalts have a continuum in titanium concentrations, with the highest concentration rocks being the least abundant.
Internal structure
The current model of the interior of the Moon was derived using seismometers left behind during the crewed Apollo program missions, as well as investigations of the Moon's gravity field and rotation.
The mass of the Moon is sufficient to eliminate any voids within the interior, so it is estimated to be composed of solid rock throughout. Its low bulk density (~3346 kg m−3) indicates a low metal abundance. Mass and moment of inertia constraints indicate that the Moon likely has an iron core that is less than about 450 km in radius. Studies of the Moon's physical librations (small perturbations to its rotation) furthermore indicate that the core is still molten. Most planetary bodies and moons have iron cores that are about half the size of the body. The Moon is thus anomalous in having a core whose size is only about one quarter of its radius.
The crust of the Moon is on average about 50 km thick (though this is uncertain by about ±15 km). It is estimated that the far-side crust is on average thicker than the near side by about 15 km. Seismology has constrained the thickness of the crust only near the Apollo 12 and Apollo 14 landing sites. Although the initial Apollo-era analyses suggested a crustal thickness of about 60 km at this site, recent reanalyses of this data suggest that it is thinner, somewhere between about 30 and 45 km.
Magnetic field
Compared with Earth, the Moon has a weak external magnetic field. Other significant differences are that the Moon does not currently have a dipolar magnetic field (as would be generated by a geodynamo in its core), and the magnetizations that are present are almost entirely crustal in origin. One hypothesis holds that the crustal magnetizations were acquired early in lunar history when a geodynamo was still operating. However, the lunar core's small size is a potential obstacle to this hypothesis. Alternatively, it is possible that transient magnetic fields could be generated during impact processes on airless bodies such as the Moon. In support of this, it has been noted that the largest crustal magnetizations appear to be located near the antipodes of the largest impact basins.
Although the Moon does not have a dipolar magnetic field like Earth's, some returned rocks have strong magnetizations. Furthermore, measurements from orbit show that some portions of the lunar surface are associated with strong magnetic fields.
| Physical sciences | Solar System | Astronomy |
2099579 | https://en.wikipedia.org/wiki/Rhynchocinetidae | Rhynchocinetidae | The family Rhynchocinetidae are a group of small, reclusive red-and-white shrimp. This family typically has an upward-hinged foldable rostrum, hence its taxon name Rhynchocinetidae, which means movable beak; this gives these shrimps their common name of hinge-beak shrimps. The family contains only two genera, Cinetorhynchus and Rhynchocinetes.
Taxonomy
Rhynchocinetidae has historically been considered to include the single genus Rhynchocinetes, which was subdivided into two sub-genera. However, in 1995, Holthuis elevated the subgenus Cinetorhynchus to full generic status based on morphology. Members of Rhynchocinetes have two acute teeth on the central carina of the carapace, a supraorbital spine and no spine on the margins of the fourth and fifth abdominal somites. Cinetorhynchus differs in having three teeth on the carapace, no supraorbital spine and a single spine each on the margins of the fourth and fifth abdominal somites.
Description
Members of the family Rhynchocinetidae are colourful crustaceans, mostly banded or patterned in red and white. The rostrum is partially or completely articulated with the carapace and can be moved in a variety of ways. The advantages of this adaptation are unclear, but in combination with the spiny rostrum, it may enable the shrimp to wedge itself into a crevice and avoid being dislodged by a predator. The first pair of chelipeds are more robust than the second pair, and the carpus of the second pair is undivided. In males, particularly in larger, older individuals, the third pair of maxillipeds and the first pair of chelipeds are exceptionally large in proportion to those of females and younger males. Both pairs of chelipeds bear a bundle of spines at the tip which when retracted form a basket-like cage.
Distribution and habitat
Members of the family Rhynchocinetidae are found in both tropical and some temperate waters, from the littoral zone down to the continental shelf, inhabiting both coral and rocky reefs.
Ecology
The biology of the family has been little studied. Cinetorhynchus rigens in Bermuda was not attracted to baited traps but on examination, its faecal pellets contained mollusc shell fragments, algae and sponge spicules. It was nocturnal, hiding by day in crevices and emerging at night to feed.
| Biology and health sciences | Shrimps and prawns | Animals |
35852419 | https://en.wikipedia.org/wiki/Aphrophoridae | Aphrophoridae | The Aphrophoridae are a family of spittlebugs belonging to the order Hemiptera. There are at least 160 genera and 990 described species in Aphrophoridae.
European genera
Aphrophora Germar 1821
Lepyronia Amyot & Serville 1843
Mesoptyelus Matsumura 1904
Neophilaenus Haupt 1935
Paraphilaenus Vilbaste 1962
Peuceptyelus Sahlberg 1871
Philaenus Stål 1864
| Biology and health sciences | Hemiptera (true bugs) | Animals |
24536042 | https://en.wikipedia.org/wiki/Feces | Feces | Feces (or faeces; : faex) are the solid or semi-solid remains of food that was not digested in the small intestine, and has been broken down by bacteria in the large intestine. Feces contain a relatively small amount of metabolic waste products such as bacterially altered bilirubin, and dead epithelial cells from the lining of the gut.
Feces are discharged through the anus or cloaca during defecation.
Feces can be used as fertilizer or soil conditioner in agriculture. They can also be burned as fuel or dried and used for construction. Some medicinal uses have been found. In the case of human feces, fecal transplants or fecal bacteriotherapy are in use. Urine and feces together are called excreta.
Characteristics
The distinctive odor of feces is due to skatole, and thiols (sulfur-containing compounds), as well as amines and carboxylic acids. Skatole is produced from tryptophan via indoleacetic acid. Decarboxylation gives skatole.
The perceived bad odor of feces has been hypothesized to be a deterrent for humans, as consuming or touching it may result in sickness or infection.
Physiology
Feces are discharged through the anus or cloaca during defecation. This process requires pressures that may reach (13.3 kPa) in humans and (60 kPa) in penguins. The forces required to expel the feces are generated through muscular contractions and a build-up of gases inside the gut, prompting the sphincter to relieve the pressure and release the feces.
Ecology
After an animal has digested eaten material, the remains of that material are discharged from its body as waste. Although it is lower in energy than the food from which it is derived, feces may retain a large amount of energy, often 50% of that of the original food. This means that of all food eaten, a significant amount of energy remains for the decomposers of ecosystems.
Many organisms feed on feces, from bacteria to fungi to insects such as dung beetles, who can sense odors from long distances. Some may specialize in feces, while others may eat other foods. Feces serve not only as a basic food, but also as a supplement to the usual diet of some animals. This process is known as coprophagia, and occurs in various animal species such as young elephants eating the feces of their mothers to gain essential gut flora, or by other animals such as dogs, rabbits, and monkeys.
Feces and urine, which reflect ultraviolet light, are important to raptors such as kestrels, who can see the near ultraviolet and thus find their prey by their middens and territorial markers.
Seeds also may be found in feces. Animals who eat fruit are known as frugivores. An advantage for a plant in having fruit is that animals will eat the fruit and unknowingly disperse the seed in doing so. This mode of seed dispersal is highly successful, as seeds dispersed around the base of a plant are unlikely to succeed and often are subject to heavy predation. Provided the seed can withstand the pathway through the digestive system, it is not only likely to be far away from the parent plant, but is even provided with its own fertilizer.
Organisms that subsist on dead organic matter or detritus are known as detritivores, and play an important role in ecosystems by recycling organic matter back into a simpler form that plants and other autotrophs may absorb once again. This cycling of matter is known as the biogeochemical cycle. To maintain nutrients in soil it is therefore important that feces returns to the area from which they came, which is not always the case in human society where food may be transported from rural areas to urban populations and then feces disposed of into a river or sea.
Human feces
Depending on the individual and the circumstances, human beings may defecate several times a day, every day, or once every two or three days. Extensive hardening of the feces that interrupts this routine for several days or more is called constipation.
The appearance of human fecal matter varies according to diet and health. Normally it is semisolid, with a mucus coating. A combination of bile and bilirubin, which comes from dead red blood cells, gives feces the typical brown color.
After the meconium, the first stool expelled, a newborn's feces contains only bile, which gives it a yellow-green color. Breast feeding babies expel soft, pale yellowish, and not quite malodorous matter; but once the baby begins to eat, and the body starts expelling bilirubin from dead red blood cells, its matter acquires the familiar brown color.
At different times in their life, human beings will expel feces of different colors and textures. A stool that passes rapidly through the intestines will look greenish; lack of bilirubin will make the stool look like clay.
Uses of animal feces
Fertilizer
The feces of animals, e.g. guano and manure, often are used as fertilizer.
Energy
Dry animal dung, such as that of camel, bison and cattle, is burned as fuel in many countries.
Animals such as the giant panda and zebra possess gut bacteria capable of producing biofuel. The bacterium in question, Brocadia anammoxidans, can be used to synthesize the rocket fuel hydrazine.
Coprolites and paleofeces
A coprolite is fossilized feces and is classified as a trace fossil. In paleontology they give evidence about the diet of an animal. They were first described by William Buckland in 1829. Prior to this, they were known as "fossil fir cones" and "bezoar stones". They serve a valuable purpose in paleontology because they provide direct evidence of the predation and diet of extinct organisms. Coprolites may range in size from a few millimetres to more than 60 centimetres.
Palaeofeces are ancient feces, often found as part of archaeological excavations or surveys. Intact paleofeces of ancient people may be found in caves in arid climates and in other locations with suitable preservation conditions. These are studied to determine the diet and health of the people who produced them through the analysis of seeds, small bones, and parasite eggs found inside. Feces may contain information about the person excreting the material as well as information about the material. They also may be analyzed chemically for more in-depth information on the individual who excreted them, using lipid analysis and ancient DNA analysis. The success rate of usable DNA extraction is relatively high in paleofeces, making it more reliable than skeletal DNA retrieval.
The reason this analysis is possible at all is due to the digestive system not being entirely efficient, in the sense that not everything that passes through the digestive system is destroyed. Not all of the surviving material is recognizable, but some of it is. Generally, this material is the best indicator archaeologists can use to determine ancient diets, as no other part of the archaeological record is so direct an indicator.
A process that preserves feces in a way that they may be analyzed later is the Maillard reaction. This reaction creates a casing of sugar that preserves the feces from the elements. To extract and analyze the information contained within, researchers generally have to freeze the feces and grind it up into powder for analysis.
Other uses
Animal dung occasionally is used as a cement to make adobe (mudbrick) huts, or even in throwing sports, especially with cow and camel dung.
Kopi luwak, or civet coffee, is coffee made from coffee beans that have been eaten and excreted by Asian palm civets (Paradoxurus hermaphroditus).
Giant pandas provide fertilizer for the world's most expensive green tea. In Malaysia, tea is made from the droppings of stick insects fed on guava leaves.
In northern Thailand, elephants are used to digest coffee beans in order to make Black Ivory coffee, which is among the world's most expensive coffees. Paper is also made from elephant dung in Thailand. Haathi Chaap is a brand of paper made from elephant dung.
Dog feces was used in the tanning process of leather during the Victorian era. Collected dog feces, known as "pure", "puer", or "pewer", were mixed with water to form a substance known as "bate", because proteolytic enzymes in the dog feces helped to relax the fibrous structure of the hide before the final stages of tanning. Dog feces collectors were known as pure finders.
Elephants, hippos, koalas and pandas are born with sterile intestines, and require bacteria obtained from eating the feces of their mothers to digest vegetation.
In India, cow dung and cow urine are major ingredients of the traditional Hindu drink Panchagavya. Politician Shankarbhai Vegad stated that they can cure cancer.
Terminology
Feces is the scientific terminology, while the term stool is also commonly used in medical contexts. Outside of scientific contexts, these terms are less common, with the most common layman's term being poop or poo. The term shit is also in common use, although it is widely considered vulgar or offensive. There are many other terms, see below.
Etymology
The word faeces is the plural of the Latin word faex meaning "dregs". In most English-language usage, there is no singular form, making the word a plurale tantum; out of various major dictionaries, only one enters variation from plural agreement.
Synonyms
"Feces" is used more in biology and medicine than in other fields (reflecting science's tradition of classical Latin and Neo-Latin)
In hunting and tracking, terms such as dung, scat, spoor, and droppings normally are used to refer to non-human animal feces
In husbandry and farming, manure is common.
Stool is a common term in reference to human feces. For example, in medicine, to diagnose the presence or absence of a medical condition, a stool sample sometimes is requested for testing purposes.
The term bowel movement(s) (with each movement a defecation event) is also common in health care.
There are many synonyms in informal registers for feces, just like there are for urine. Many are euphemistic, colloquial, or both; some are profane (such as shit), whereas most belong chiefly to child-directed speech (such as poo or the palindromic word poop) or to crude humor (such as crap, dump, load and turd.).
Feces of animals
The feces of animals often have special names (some of them are slang), for example:
Non-human animals
As bulk material – dung
Individually – droppings
Cattle
Bulk material – cow dung
Individual droppings – cow pats, meadow muffins, etc.
Deer (and formerly other quarry animals) – fewmets
Wild carnivores – scat
Otter – spraint
Birds (individual) – droppings (also include urine as white crystals of uric acid)
Seabirds or bats (large accumulations) – guano
Herbivorous insects, such as caterpillars and leaf beetles – frass
Earthworms, lugworms etc. – worm castings (feces extruded at ground surface)
Feces when used as fertilizer (usually mixed with animal bedding and urine) – manure
Horses – horse manure, roadapple (before motor vehicles became common, horse droppings were a big part of the rubbish communities needed to clean off roads)
Society and culture
Feelings of disgust
In all human cultures, feces elicit varying degrees of disgust in adults. Children under two years typically have no disgust response to it, suggesting it is culturally derived. Disgust toward feces appears to be strongest in cultures where flush toilets make olfactory contact with human feces minimal. Disgust is experienced primarily in relation to the sense of taste (either perceived or imagined) and, secondarily to anything that causes a similar feeling by sense of smell, touch, or vision.
Social media
There is a Pile of Poo emoji represented in Unicode as , called or unchi-kun in Japan.
Jokes
Poop is the center of toilet humor, and is commonly an interest of young children and teenagers.
| Biology and health sciences | Gastrointestinal tract | Biology |
24536543 | https://en.wikipedia.org/wiki/Eukaryote | Eukaryote | The eukaryotes ( ) constitute the domain of Eukaryota or Eukarya, organisms whose cells have a membrane-bound nucleus. All animals, plants, fungi, seaweeds, and many unicellular organisms are eukaryotes. They constitute a major group of life forms alongside the two groups of prokaryotes: the Bacteria and the Archaea. Eukaryotes represent a small minority of the number of organisms, but given their generally much larger size, their collective global biomass is much larger than that of prokaryotes.
The eukaryotes seemingly emerged within the Asgard archaea, and are closely related to the Heimdallarchaeia. This implies that there are only two domains of life, Bacteria and Archaea, with eukaryotes incorporated among the Archaea. Eukaryotes first emerged during the Paleoproterozoic, likely as flagellated cells. The leading evolutionary theory is they were created by symbiogenesis between an anaerobic Asgard archaean and an aerobic proteobacterium, which formed the mitochondria. A second episode of symbiogenesis with a cyanobacterium created the plants, with chloroplasts.
Eukaryotic cells contain membrane-bound organelles such as the nucleus, the endoplasmic reticulum, and the Golgi apparatus. Eukaryotes may be either unicellular or multicellular. In comparison, prokaryotes are typically unicellular. Unicellular eukaryotes are sometimes called protists. Eukaryotes can reproduce both asexually through mitosis and sexually through meiosis and gamete fusion (fertilization).
Etymology
The word 'Eukaryote' is derived from the Greek words "eu" (εὖ) meaning "true" or "good" and "karyon" (κάρυον) meaning "nut" or "kernel," referring to the nucleus of a cell.
Diversity
Eukaryotes are organisms that range from microscopic single cells, such as picozoans under 3 micrometres across, to animals like the blue whale, weighing up to 190 tonnes and measuring up to long, or plants like the coast redwood, up to tall. Many eukaryotes are unicellular; the informal grouping called protists includes many of these, with some multicellular forms like the giant kelp up to long. The multicellular eukaryotes include the animals, plants, and fungi, but again, these groups too contain many unicellular species. Eukaryotic cells are typically much larger than those of prokaryotes—the bacteria and the archaea—having a volume of around 10,000 times greater. Eukaryotes represent a small minority of the number of organisms, but, as many of them are much larger, their collective global biomass (468 gigatons) is far larger than that of prokaryotes (77 gigatons), with plants alone accounting for over 81% of the total biomass of Earth.
The eukaryotes are a diverse lineage, consisting mainly of microscopic organisms. Multicellularity in some form has evolved independently at least 25 times within the eukaryotes. Complex multicellular organisms, not counting the aggregation of amoebae to form slime molds, have evolved within only six eukaryotic lineages: animals, symbiomycotan fungi, brown algae, red algae, green algae, and land plants. Eukaryotes are grouped by genomic similarities, so that groups often lack visible shared characteristics.
Distinguishing features
Nucleus
The defining feature of eukaryotes is that their cells have a well-defined, membrane-bound nuclei, distinguishing them from prokaryotes that lack such a structure. Eukaryotic cells have a variety of internal membrane-bound structures, called organelles, and a cytoskeleton which defines the cell's organization and shape. The nucleus stores the cell's DNA, which is divided into linear bundles called chromosomes; these are separated into two matching sets by a microtubular spindle during nuclear division, in the distinctively eukaryotic process of mitosis.
Biochemistry
Eukaryotes differ from prokaryotes in multiple ways, with unique biochemical pathways such as sterane synthesis. The eukaryotic signature proteins have no homology to proteins in other domains of life, but appear to be universal among eukaryotes. They include the proteins of the cytoskeleton, the complex transcription machinery, the membrane-sorting systems, the nuclear pore, and some enzymes in the biochemical pathways.
Internal membranes
Eukaryote cells include a variety of membrane-bound structures, together forming the endomembrane system. Simple compartments, called vesicles and vacuoles, can form by budding off other membranes. Many cells ingest food and other materials through a process of endocytosis, where the outer membrane invaginates and then pinches off to form a vesicle. Some cell products can leave in a vesicle through exocytosis.
The nucleus is surrounded by a double membrane known as the nuclear envelope, with nuclear pores that allow material to move in and out. Various tube- and sheet-like extensions of the nuclear membrane form the endoplasmic reticulum, which is involved in protein transport and maturation. It includes the rough endoplasmic reticulum, covered in ribosomes which synthesize proteins; these enter the interior space or lumen. Subsequently, they generally enter vesicles, which bud off from the smooth endoplasmic reticulum. In most eukaryotes, these protein-carrying vesicles are released and further modified in stacks of flattened vesicles (cisternae), the Golgi apparatus.
Vesicles may be specialized; for instance, lysosomes contain digestive enzymes that break down biomolecules in the cytoplasm.
Mitochondria
Mitochondria are organelles in eukaryotic cells. The mitochondrion is commonly called "the powerhouse of the cell", for its function providing energy by oxidising sugars or fats to produce the energy-storing molecule ATP. Mitochondria have two surrounding membranes, each a phospholipid bilayer, the inner of which is folded into invaginations called cristae where aerobic respiration takes place.
Mitochondria contain their own DNA, which has close structural similarities to bacterial DNA, from which it originated, and which encodes rRNA and tRNA genes that produce RNA which is closer in structure to bacterial RNA than to eukaryote RNA.
Some eukaryotes, such as the metamonads Giardia and Trichomonas, and the amoebozoan Pelomyxa, appear to lack mitochondria, but all contain mitochondrion-derived organelles, like hydrogenosomes or mitosomes, having lost their mitochondria secondarily. They obtain energy by enzymatic action in the cytoplasm. It is thought that mitochondria developed from prokaryotic cells which became endosymbionts living inside eukaryotes.
Plastids
Plants and various groups of algae have plastids as well as mitochondria. Plastids, like mitochondria, have their own DNA and are developed from endosymbionts, in this case cyanobacteria. They usually take the form of chloroplasts which, like cyanobacteria, contain chlorophyll and produce organic compounds (such as glucose) through photosynthesis. Others are involved in storing food. Although plastids probably had a single origin, not all plastid-containing groups are closely related. Instead, some eukaryotes have obtained them from others through secondary endosymbiosis or ingestion. The capture and sequestering of photosynthetic cells and chloroplasts, kleptoplasty, occurs in many types of modern eukaryotic organisms.
Cytoskeletal structures
The cytoskeleton provides stiffening structure and points of
attachment for motor structures that enable the cell to move, change shape, or transport materials. The motor structures are microfilaments of actin and actin-binding proteins, including α-actinin, fimbrin, and filamin are present in submembranous cortical layers and bundles. Motor proteins of microtubules, dynein and kinesin, and myosin of actin filaments, provide dynamic character of the network.
Many eukaryotes have long slender motile cytoplasmic projections, called flagella, or multiple shorter structures called cilia. These organelles are variously involved in movement, feeding, and sensation. They are composed mainly of tubulin, and are entirely distinct from prokaryotic flagella. They are supported by a bundle of microtubules arising from a centriole, characteristically arranged as nine doublets surrounding two singlets. Flagella may have hairs (mastigonemes), as in many stramenopiles. Their interior is continuous with the cell's cytoplasm.
Centrioles are often present, even in cells and groups that do not have flagella, but conifers and flowering plants have neither. They generally occur in groups that give rise to various microtubular roots. These form a primary component of the cytoskeleton, and are often assembled over the course of several cell divisions, with one flagellum retained from the parent and the other derived from it. Centrioles produce the spindle during nuclear division.
Cell wall
The cells of plants, algae, fungi and most chromalveolates, but not animals, are surrounded by a cell wall. This is a layer outside the cell membrane, providing the cell with structural support, protection, and a filtering mechanism. The cell wall also prevents over-expansion when water enters the cell.
The major polysaccharides making up the primary cell wall of land plants are cellulose, hemicellulose, and pectin. The cellulose microfibrils are linked together with hemicellulose, embedded in a pectin matrix. The most common hemicellulose in the primary cell wall is xyloglucan.
Sexual reproduction
Eukaryotes have a life cycle that involves sexual reproduction, alternating between a haploid phase, where only one copy of each chromosome is present in each cell, and a diploid phase, with two copies of each chromosome in each cell. The diploid phase is formed by fusion of two haploid gametes, such as eggs and spermatozoa, to form a zygote; this may grow into a body, with its cells dividing by mitosis, and at some stage produce haploid gametes through meiosis, a division that reduces the number of chromosomes and creates genetic variability. There is considerable variation in this pattern. Plants have both haploid and diploid multicellular phases. Eukaryotes have lower metabolic rates and longer generation times than prokaryotes, because they are larger and therefore have a smaller surface area to volume ratio.
The evolution of sexual reproduction may be a primordial characteristic of eukaryotes. Based on a phylogenetic analysis, Dacks and Roger have proposed that facultative sex was present in the group's common ancestor. A core set of genes that function in meiosis is present in both Trichomonas vaginalis and Giardia intestinalis, two organisms previously thought to be asexual. Since these two species are descendants of lineages that diverged early from the eukaryotic evolutionary tree, core meiotic genes, and hence sex, were likely present in the common ancestor of eukaryotes. Species once thought to be asexual, such as Leishmania parasites, have a sexual cycle. Amoebae, previously regarded as asexual, may be anciently sexual; while present-day asexual groups could have arisen recently.
Evolution
History of classification
In antiquity, the two lineages of animals and plants were recognized by Aristotle and Theophrastus. The lineages were given the taxonomic rank of kingdom by Linnaeus in the 18th century. Though he included the fungi with plants with some reservations, it was later realized that they are quite distinct and warrant a separate kingdom. The various single-cell eukaryotes were originally placed with plants or animals when they became known. In 1818, the German biologist Georg A. Goldfuss coined the word Protozoa to refer to organisms such as ciliates, and this group was expanded until Ernst Haeckel made it a kingdom encompassing all single-celled eukaryotes, the Protista, in 1866. The eukaryotes thus came to be seen as four kingdoms:
Kingdom Protista
Kingdom Plantae
Kingdom Fungi
Kingdom Animalia
The protists were at that time thought to be "primitive forms", and thus an evolutionary grade, united by their primitive unicellular nature. Understanding of the oldest branchings in the tree of life only developed substantially with DNA sequencing, leading to a system of domains rather than kingdoms as top level rank being put forward by Carl Woese, Otto Kandler, and Mark Wheelis in 1990, uniting all the eukaryote kingdoms in the domain "Eucarya", stating, however, that eukaryotes' will continue to be an acceptable common synonym". In 1996, the evolutionary biologist Lynn Margulis proposed to replace kingdoms and domains with "inclusive" names to create a "symbiosis-based phylogeny", giving the description "Eukarya (symbiosis-derived nucleated organisms)".
Phylogeny
By 2014, a rough consensus started to emerge from the phylogenomic studies of the previous two decades. The majority of eukaryotes can be placed in one of two large clades dubbed Amorphea (similar in composition to the unikont hypothesis) and the Diphoda (formerly bikonts), which includes plants and most algal lineages. A third major grouping, the Excavata, has been abandoned as a formal group as it is paraphyletic. The proposed phylogeny below includes only one group of excavates (Discoba), and incorporates the 2021 proposal that picozoans are close relatives of rhodophytes. The Provora are a group of microbial predators discovered in 2022.
Origin of eukaryotes
The origin of the eukaryotic cell, or eukaryogenesis, is a milestone in the evolution of life, since eukaryotes include all complex cells and almost all multicellular organisms. The last eukaryotic common ancestor (LECA) is the hypothetical origin of all living eukaryotes, and was most likely a biological population, not a single individual. The LECA is believed to have been a protist with a nucleus, at least one centriole and flagellum, facultatively aerobic mitochondria, sex (meiosis and syngamy), a dormant cyst with a cell wall of chitin or cellulose, and peroxisomes.
An endosymbiotic union between a motile anaerobic archaean and an aerobic alphaproteobacterium gave rise to the LECA and all eukaryotes, with mitochondria. A second, much later endosymbiosis with a cyanobacterium gave rise to the ancestor of plants, with chloroplasts.
The presence of eukaryotic biomarkers in archaea points towards an archaeal origin. The genomes of Asgard archaea have plenty of eukaryotic signature protein genes, which play a crucial role in the development of the cytoskeleton and complex cellular structures characteristic of eukaryotes. In 2022, cryo-electron tomography demonstrated that Asgard archaea have a complex actin-based cytoskeleton, providing the first direct visual evidence of the archaeal ancestry of eukaryotes.
Fossils
The timing of the origin of eukaryotes is hard to determine, but the discovery of Qingshania magnificia, the earliest multicelluar eukaryote from North China which lived 1.635 billion years ago, suggests that the crown group eukaryotes originated from the late Paleoproterozoic (Statherian). The earliest unequivocal unicellular eukaryotes, Tappania plana, Shuiyousphaeridium macroreticulatum, Dictyosphaera macroreticulata, Germinosphaera alveolata, and Valeria lophostriata from North China, lived approximately 1.65 billion years ago.
Some acritarchs are known from at least 1.65 billion years ago, and a fossil, Grypania, which may be an alga, is as much as 2.1 billion years old. The "problematic" fossil Diskagma has been found in paleosols 2.2 billion years old.
The Neoarchean fossil Thuchomyces shares similarities with eukaryotes, specifically fungi. It especially resembles the problematic fossil Diskagma, with hyphae and multiple differentiated layers. However, it is over 600 million years older than all other possible eukaryotes, and many of its "eukaryote features" are not specific to the clade, meaning it is almost certainly a microbial mat instead.
Structures proposed to represent "large colonial organisms" have been found in the black shales of the Palaeoproterozoic such as the Francevillian B Formation, in Gabon, dubbed the "Francevillian biota" which is dated at 2.1 billion years old. However, the status of these structures as fossils is contested, with other authors suggesting that they might represent pseudofossils. The oldest fossils than can unambiguously be assigned to eukaryotes are from the Ruyang Group of China, dating to approximately 1.8-1.6 billion years ago. Fossils that are clearly related to modern groups start appearing an estimated 1.2 billion years ago, in the form of red algae, though recent work suggests the existence of fossilized filamentous algae in the Vindhya basin dating back perhaps to 1.6 to 1.7 billion years ago.
The presence of steranes, eukaryotic-specific biomarkers, in Australian shales previously indicated that eukaryotes were present in these rocks dated at 2.7 billion years old, but these Archaean biomarkers have been rebutted as later contaminants. The oldest valid biomarker records are only around 800 million years old. In contrast, a molecular clock analysis suggests the emergence of sterol biosynthesis as early as 2.3 billion years ago. The nature of steranes as eukaryotic biomarkers is further complicated by the production of sterols by some bacteria.
Whenever their origins, eukaryotes may not have become ecologically dominant until much later; a massive increase in the zinc composition of marine sediments has been attributed to the rise of substantial populations of eukaryotes, which preferentially consume and incorporate zinc relative to prokaryotes, approximately a billion years after their origin (at the latest).
| Biology and health sciences | Biology | null |
24539067 | https://en.wikipedia.org/wiki/Ardi | Ardi | Ardi (ARA-VP-6/500) is the designation of the fossilized skeletal remains of an Ardipithecus ramidus, thought to be an early human-like female anthropoid 4.4 million years old. It is the most complete early hominid specimen, with most of the skull, teeth, pelvis, hands and feet, more complete than the previously known Australopithecus afarensis specimen called "Lucy". In all, 125 different pieces of fossilized bone were found.
Discovery
The Ardi skeleton was discovered at Aramis in the arid badlands near the Awash River in Ethiopia in 1994 by a college student, Yohannes Haile-Selassie, when he uncovered a partial piece of a hand bone. The discovery was made by a team of scientists led by UC Berkeley anthropologist, Tim D. White, and was analyzed by an international group of scientists that included Owen Lovejoy heading the biology team. On 1 October 2009, the journal Science published an open-access collection of eleven articles, detailing many aspects of A. ramidus and its environment. Her fossils were also found near animal remains which indicated that she inhabited a forest type of environment, contrary to the theory that bipedalism originated in savannahs.
Ardi was not the first fossil of A. ramidus to come to light. The first ones were found in Ethiopia in 1992, but it took 17 years to assess their significance.
Etymology
The word Ardi means "ground floor" and the word ramid means "root" in the Afar language, suggesting that Ardi lived on the ground and was the root of the family tree of humanity.
Description
Ardi weighed about , and could be up to tall. Although she is a biped, Ardi had both opposable big toes and thumbs in order to climb trees. It is speculated that her bipedality impeded movement, but enabled her to bear more offspring.
Although it is not yet clear how Ardi's species is related to Homo sapiens, the discovery is of great significance and added much to the debate on Ardipithecus and its place in human evolution. With regards to Ardi's body composition, archaeologists note that she is unique in that she possesses traits that are characteristic of both extinct primates and early hominids. It is still a point of debate whether Ardi was capable of bipedal movement. Ardi's divergent big toes are not characteristic of a biped. However, the found remains of her legs, feet, pelvis, and hands suggested that she walked upright when on the ground but was a quadruped when moving around trees. Her big toe, for example, spreads out quite a bit from her foot to better grasp tree limbs. Unlike chimpanzees, however, her foot contains a unique small bone inside a tendon which kept the big toe stronger. When seen along with Ardi's other bone structures, this unique bone would have helped her walk bipedally, though less efficiently than Lucy. Her wrist bones also provided her with flexibility but the palm bones were short. This suggests that Ardi did not walk on her knuckles and only used her palms to move along tree branches.
Some of Ardi's teeth are still connected to her jawbone and show enamel wear suggesting a diet consisting of fruit and nuts. The canine teeth of A. ramidus are smaller, and equal in size between males and females. This suggests reduced male-to-male conflict, pair-bonding, and increased parental investment. "Thus fundamental reproductive and social behavioral changes probably occurred in hominids long before they had enlarged brains and began to use stone tools."
Pelvis
Although Ardipithecus had more ancestral hands, feet, and limbs, Ardi's pelvis gives a different perspective. The parts of Ardi's pelvis that were recovered include her left hip, her right ilium, and a fragment of her distal sacrum. A shorter ilium and a curve in the lower spine were the characteristics gathered from these partial remains that indicate Ardi, and the Ardipithecus ramidus species, had the ability to walk upright. The shift to bipedality is only beginning to emerge in Ardi because there are characteristics in Ardi's pelvis that are both found in all later hominids and characteristics that are found in extant African apes. A characteristic that is found in Ardi and in all later hominids is a separate growth site for the anterior inferior iliac spine. A similar ischial structure is a characteristic found in Ardi and in extant African apes. This mixture of characteristics indicates Ardi's bipedality was an earlier version of bipedalism compared to later hominids like Lucy. Regardless of how ancestral Ardi's bipedality was, these characteristics found in Ardi's pelvis show bipedalism was well underway by 4.4 million years ago, even with the ability for arboreal locomotion still present in the hands and limbs.
Foot
Ardi's foot is a special area of interest when examining the evolution of bipedalism in early Hominids, and the bipedality of Ardipithecus ramidus, because all five toes do not line up. The remains of the foot from Ardi and other Ardipithecus ramidus specimens that can be studied includes "a talus, medial and intermediate cuneiforms, cuboid, first, second, third, and fifth metatarsals, and several phalanges." The foot of Ardi contains an opposable hallux (big toe) that is similar to chimpanzees. This opposable hallux is believed to have been used to aid in tree climbing. On the outside, Ardi's foot may look like it belongs with other Apes, but on the inside, Ardi's foot contains a bone called the os peroneum, which allows the bottom of the foot to be more rigid. The rigidity of the bottom of the foot was believed to allow Ardi to walk upright, and the other four toes that were aligned performed the "toe off" action during a bipedal motion. The combination of features found in Ardi's and other Ardipithecus ramidus foot bones captures a moment in time where these primitive primates were beginning to leave the trees and spending longer periods of time on the ground.
| Biology and health sciences | Australopithecines | Biology |
21556842 | https://en.wikipedia.org/wiki/Photographic%20film | Photographic film | Photographic film is a strip or sheet of transparent film base coated on one side with a gelatin emulsion containing microscopically small light-sensitive silver halide crystals. The sizes and other characteristics of the crystals determine the sensitivity, contrast, and resolution of the film. Film is typically segmented in frames, that give rise to separate photographs.
The emulsion will gradually darken if left exposed to light, but the process is too slow and incomplete to be of any practical use. Instead, a very short exposure to the image formed by a camera lens is used to produce only a very slight chemical change, proportional to the amount of light absorbed by each crystal. This creates an invisible latent image in the emulsion, which can be chemically developed into a visible photograph. In addition to visible light, all films are sensitive to ultraviolet light, X-rays, gamma rays, and high-energy particles. Unmodified silver halide crystals are sensitive only to the blue part of the visible spectrum, producing unnatural-looking renditions of some colored subjects. This problem was resolved with the discovery that certain dyes, called sensitizing dyes, when adsorbed onto the silver halide crystals made them respond to other colors as well. First orthochromatic (sensitive to blue and green) and finally panchromatic (sensitive to all visible colors) films were developed. Panchromatic film renders all colors in shades of gray approximately matching their subjective brightness. By similar techniques, special-purpose films can be made sensitive to the infrared (IR) region of the spectrum.
In black-and-white photographic film, there is usually one layer of silver halide crystals. When the exposed silver halide grains are developed, the silver halide crystals are converted to metallic silver, which blocks light and appears as the black part of the film negative. Color film has at least three sensitive layers, incorporating different combinations of sensitizing dyes. Typically the blue-sensitive layer is on top, followed by a yellow filter layer to stop any remaining blue light from affecting the layers below. Next comes a green-and-blue sensitive layer, and a red-and-blue sensitive layer, which record the green and red images respectively. During development, the exposed silver halide crystals are converted to metallic silver, just as with black-and-white film. But in a color film, the by-products of the development reaction simultaneously combine with chemicals known as color couplers that are included either in the film itself or in the developer solution to form colored dyes. Because the by-products are created in direct proportion to the amount of exposure and development, the dye clouds formed are also in proportion to the exposure and development. Following development, the silver is converted back to silver halide crystals in the bleach step. It is removed from the film during the process of fixing the image on the film with a solution of ammonium thiosulfate or sodium thiosulfate (hypo or fixer). Fixing leaves behind only the formed color dyes, which combine to make up the colored visible image. Later color films, like Kodacolor II, have as many as 12 emulsion layers, with upwards of 20 different chemicals in each layer.
Photographic film and film stock tend to be similar in composition and speed, but often not in other parameters such as frame size and length. Silver halide photographic paper is also similar to photographic film.
Before the emergence of digital photography, photographs on film had to be developed to produce negatives or projectable slides, and negatives had to be printed as positive images, usually in enlarged form. This was usually done by photographic laboratories, but many amateurs did their own processing.
Characteristics of film
Film basics
There are several types of photographic film, including:
Print film, when developed, yields transparent negatives with the light and dark areas and colors (if color film is used) inverted to their respective complementary colors. This type of film is designed to be printed onto photographic paper, usually by means of an enlarger but in some cases by contact printing. The paper is then itself developed. The second inversion that results restores light, shade and color to their normal appearance. Color negatives incorporate an orange color correction mask that compensates for unwanted dye absorptions and improves color accuracy in the prints. Although color processing is more complex and temperature-sensitive than black-and-white processing, the wide availability of commercial color processing and scarcity of service for black-and-white prompted the design of some black-and-white films which are processed in exactly the same way as standard color film.
Color reversal film produces positive transparencies, also known as diapositives. Transparencies can be reviewed with the aid of a magnifying loupe and a lightbox. If mounted in small metal, plastic or cardboard frames for use in a slide projector or slide viewer they are commonly called slides. Reversal film is often marketed as "slide film". Large-format color reversal sheet film is used by some professional photographers, typically to originate very-high-resolution imagery for digital scanning into color separations for mass photomechanical reproduction. Photographic prints can be produced from reversal film transparencies, but positive-to-positive print materials for doing this directly (e.g. Ektachrome paper, Cibachrome/Ilfochrome) have all been discontinued, so it now requires the use of an internegative to convert the positive transparency image into a negative transparency, which is then printed as a positive print.
Black-and-white reversal film exists but is very uncommon. Conventional black-and-white negative film can be reversal-processed to produce black-and-white slides, as by dr5 Chrome. Although kits of chemicals for black-and-white reversal processing may no longer be available to amateur darkroom enthusiasts, an acid bleaching solution, the only unusual component which is essential, is easily prepared from scratch. Black-and-white transparencies may also be produced by printing negatives onto special positive print film, still available from some specialty photographic supply dealers.
In order to produce a usable image, the film needs to be exposed properly. The amount of exposure variation that a given film can tolerate, while still producing an acceptable level of quality, is called its exposure latitude. Color print film generally has greater exposure latitude than other types of film. Additionally, because print film must be printed to be viewed, after-the-fact corrections for imperfect exposure are possible during the printing process.
The concentration of dyes or silver halide crystals remaining on the film after development is referred to as optical density, or simply density; the optical density is proportional to the logarithm of the optical transmission coefficient of the developed film. A dark image on the negative is of higher density than a more transparent image.
Most films are affected by the physics of silver grain activation (which sets a minimum amount of light required to expose a single grain) and by the statistics of random grain activation by photons. The film requires a minimum amount of light before it begins to expose, and then responds by progressive darkening over a wide dynamic range of exposure until all of the grains are exposed, and the film achieves (after development) its maximum optical density.
Over the active dynamic range of most films, the density of the developed film is proportional to the logarithm of the total amount of light to which the film was exposed, so the transmission coefficient of the developed film is proportional to a power of the reciprocal of the brightness of the original exposure. The plot of the density of the film image against the log of the exposure is known as an H&D curve. This effect is due to the statistics of grain activation: as the film becomes progressively more exposed, each incident photon is less likely to impact a still-unexposed grain, yielding the logarithmic behavior. A simple, idealized statistical model yields the equation , where light is proportional to the number of photons hitting a unit area of film, k is the probability of a single photon striking a grain (based on the size of the grains and how closely spaced they are), and density is the proportion of grains that have been hit by at least one photon. The relationship between density and log exposure is linear for photographic films except at the extreme ranges of maximum exposure (D-max) and minimum exposure (D-min) on an H&D curve, so the curve is characteristically S-shaped (as opposed to digital camera sensors which have a linear response through the effective exposure range). The sensitivity (i.e., the ISO speed) of a film can be affected by changing the length or temperature of development, which would move the H&D curve to the left or right (see figure).
If parts of the image are exposed heavily enough to approach the maximum density possible for a print film, then they will begin losing the ability to show tonal variations in the final print. Usually those areas will be considered overexposed and will appear as featureless white on the print. Some subject matter is tolerant of very heavy exposure. For example, sources of brilliant light, such as a light bulb or the sun, generally appear best as a featureless white on the print.
Likewise, if part of an image receives less than the beginning threshold level of exposure, which depends upon the film's sensitivity to lightor speedthe film there will have no appreciable image density, and will appear on the print as a featureless black. Some photographers use their knowledge of these limits to determine the optimum exposure for a photograph; for one example, see the Zone System. Most automatic cameras instead try to achieve a particular average density.
Color films can have many layers. The film base can have an antihalation layer applied to it or be dyed. This layer prevents light from reflecting from within the film, increasing image quality. This also can make films exposable on only one side, as it prevents exposure from behind the film. This layer is bleached after development to make it clear, thus making the film transparent. The antihalation layer, besides having a black colloidal silver sol pigment for absorbing light, can also have two UV absorbents to improve lightfastness of the developed image, an oxidized developer scavenger, dyes for compensating for optical density during printing, solvents, gelatin and disodium salt of 3,5-
disulfocatechol. If applied to the back of the film, it also serves to prevent scratching, as an antistatic measure due to its conductive carbon content, and as a lubricant to help transport the film through mechanisms. The antistatic property is necessary to prevent the film from getting fogged under low humidity, and mechanisms to avoid static are present in most if not all films. If applied on the back it is removed during film processing. If applied it may be on the back of the film base in triacetate film bases or in the front in PET film bases, below the emulsion stack. An anticurl layer and a separate antistatic layer may be present in thin high resolution films that have the antihalation layer below the emulsion. PET film bases are often dyed, specially because PET can serve as a light pipe; black and white film bases tend to have a higher level of dying applied to them. The film base needs to be transparent but with some density, perfectly flat, insensitive to light, chemically stable, resistant to tearing and strong enough to be handled manually and by camera mechanisms and film processing equipment, while being chemically resistant to moisture and the chemicals used during processing without losing strength, flexibility or changing in size.
The subbing layer is essentially an adhesive that allows the subsequent layers to stick to the film base. The film base was initially made of highly flammable cellulose nitrate, which was replaced by cellulose acetate films, often cellulose triacetate film (safety film), which in turn was replaced in many films (such as all print films, most duplication films and some other specialty films) by a PET (polyethylene terephthalate) plastic film base. Films with a triacetate base can suffer from vinegar syndrome, a decomposition process accelerated by warm and humid conditions, that releases acetic acid which is the characteristic component of vinegar, imparting the film a strong vinegar smell, accelerating damage within the film and possibly even damaging surrounding metal and films. Films are usually spliced using a special adhesive tape; those with PET layers can be ultrasonically spliced or their ends melted and then spliced.
The emulsion layers of films are made by dissolving pure silver in nitric acid to form silver nitrate crystals, which are mixed with other chemicals to form silver halide grains, which are then suspended in gelatin and applied to the film base. The size and hence the light sensitivity of these grains determines the speed of the film; since films contain real silver (as silver halide), faster films with larger crystals are more expensive and potentially subject to variations in the price of silver metal. Also, faster films have more grain, since the grains (crystals) are larger. Each crystal is often 0.2 to 2 microns in size; in color films, the dye clouds that form around the silver halide crystals are often 25 microns across. The crystals can be shaped as cubes, flat rectangles, tetradecadedra, or be flat and resemble a triangle with or without clipped edges; this type of crystal is known as a T-grain crystal or a tabular grain (T-grains). Films using T-grains are more sensitive to light without using more silver halide since they increase the surface area exposed to light by making the crystals flatter and larger in footprint instead of simply increasing their volume.
T-grains can also have a hexagonal shape. These grains also have reduced sensitivity to blue light which is an advantage since silver halide is most sensitive to blue light than other colors of light. This was traditionally solved by the addition of a blue-blocking filter layer in the film emulsion, but T-grains have allowed this layer to be removed. Also the grains may have a "core" and "shell" where the core, made of silver iodobromide, has higher iodine content than the shell, which improves light sensitivity, these grains are known as Σ-Grains.
The exact silver halide used is either silver bromide or silver bromochloroiodide, or a combination of silver bromide, chloride and iodide. Silver iodobromide may be used as a silver halide.
Silver halide crystals can be made in several shapes for use in photographic films. For example, AgBrCl hexagonal tabular grains can be used for color negative films, AgBr octahedral grains can be used for instant color photography films, AgBrl cubo-octahedral grains can be used for color reversal films, AgBr hexagonal tabular grains can be used for medical X-ray films, and AgBrCl cubic grains can be used for graphic arts films.
In color films, each emulsion layer has silver halide crystals that are sensitized to one particular color (wavelength of light) vía sentizing dyes, to that they will be made sensitive to only one color of light, and not to others, since silver halide particles are intrinsically sensitive only to wavelengths below 450 nm (which is blue light). The sensitizing dyes are absorbed at dislocations in the silver halide particles in the emulsion on the film. The sensitizing dyes may be supersensitized with a supersensitizing dye, that assists the function of the sensitizing dye and improves the efficiency of photon capture by silver halide. Each layer has a different type of color dye forming coupler: in the blue sensitive layer, the coupler forms a yellow dye; in the green sensitive layer the coupler forms a magenta dye, and in the red sensitive layer the coupler forms a cyan dye. Color films often have an UV blocking layer. Each emulsion layer in a color film may itself have three layers: a slow, medium and fast layer, to allow the film to capture higher contrast images.
The color dye couplers are inside oil droplets dispersed in the emulsion around silver halide crystals, forming a silver halide grain. Here the oil droplets act as a surfactant, also protecting the couplers from chemical reactions with the silver halide and from the surrounding gelatin. During development, oxidized developer diffuses into the oil droplets and combines with the dye couplers to form dye clouds; the dye clouds only form around unexposed silver halide crystals. The fixer then removes the silver halide crystals leaving only the dye clouds: this means that developed color films may not contain silver while undeveloped films do contain silver; this also means that the fixer can start to contain silver which can then be removed through electrolysis. Color films also contain light filters to filter out certain colors as the light passes through the film: often there is a blue light filter between the blue and green sensitive layers and a yellow filter before the red sensitive layer; in this way each layer is made sensitive to only a certain color of light.
The couplers need to be made resistant to diffusion (non-diffusible) so that they will not move between the layers of the film and thus cause incorrect color rendition as the couplers are specific to either cyan, magenta or yellow colors. This is done by making couplers with a ballast group such as a lipophilic group (oil-protected) and applying them in oil droplets to the film, or a hydrophilic group, or in a polymer layer such as a loadable latex layer with oil-protected couplers, in which case they are considered to be polymer-protected.
The color couplers may be colorless and be chromogenic or be colored. Colored couplers are used to improve the color reproduction of film. The first coupler which is used in the blue layer remains colorless to allow all light to pass through, but the coupler used in the green layer is colored yellow, and the coupler used in the red layer is light pink. Yellow was chosen to block any remaining blue light from exposing the underlying green and red layers (since yellow can be made from green and red). Each layer should only be sensitive to a single color of light and allow all others to pass through. Because of these colored couplers, the developed film appears orange. Colored couplers mean that corrections through color filters need to be applied to the image before printing. Printing can be carried out by using an optical enlarger, or by scanning the image, correcting it using software and printing it using a digital printer.
Kodachrome films have no couplers; the dyes are instead formed by a long sequence of steps, limiting adoption among smaller film processing companies.
Black and white films are very simple by comparison, only consisting of silver halide crystals suspended in a gelatin emulsion which sits on a film base with an antihalation back.
Many films contain a top supercoat layer to protect the emulsion layers from damage. Some manufacturers manufacture their films with daylight, tungsten (named after the tungsten filament of incandescent and halogen lamps) or fluorescent lighting in mind, recommending the use of lens filters, light meters and test shots in some situations to maintain color balance, or by recommending the division of the ISO value of the film by the distance of the subject from the camera to get an appropriate f-number value to be set in the lens.
Examples of Color films are Kodachrome, often processed using the K-14 process, Kodacolor, Ektachrome, which is often processed using the E-6 process and Fujifilm Superia, which is processed using the C-41 process. The chemicals and the color dye couplers on the film may vary depending on the process used to develop the film.
Film speed
Film speed describes a film's threshold sensitivity to light. The international standard for rating film speed is the ISO scale, which combines both the ASA speed and the DIN speed in the format ASA/DIN. Using ISO convention film with an ASA speed of 400 would be labeled 400/27°. A fourth naming standard is GOST, developed by the Russian standards authority. See the film speed article for a table of conversions between ASA, DIN, and GOST film speeds.
Common film speeds include ISO 25, 50, 64, 100, 160, 200, 400, 800 and 1600. Consumer print films are usually in the ISO 100 to ISO 800 range. Some films, like Kodak's Technical Pan, are not ISO rated and therefore careful examination of the film's properties must be made by the photographer before exposure and development. ISO 25 film is very "slow", as it requires much more exposure to produce a usable image than "fast" ISO 800 film. Films of ISO 800 and greater are thus better suited to low-light situations and action shots (where the short exposure time limits the total light received). The benefit of slower film is that it usually has finer grain and better color rendition than fast film. Professional photographers of static subjects such as portraits or landscapes usually seek these qualities, and therefore require a tripod to stabilize the camera for a longer exposure. A professional photographing subjects such as rapidly moving sports or in low-light conditions will inevitably choose a faster film.
A film with a particular ISO rating can be push-processed, or "pushed", to behave like a film with a higher ISO, by developing for a longer amount of time or at a higher temperature than usual. More rarely, a film can be "pulled" to behave like a "slower" film. Pushing generally coarsens grain and increases contrast, reducing dynamic range, to the detriment of overall quality. Nevertheless, it can be a useful tradeoff in difficult shooting environments, if the alternative is no usable shot at all.
Special films
Instant photography, as popularized by Polaroid, uses a special type of camera and film that automates and integrates development, without the need of further equipment or chemicals. This process is carried out immediately after exposure, as opposed to regular film, which is developed afterwards and requires additional chemicals. See instant film.
Films can be made to record non-visible ultraviolet (UV) and infrared (IR) radiation. These films generally require special equipment; for example, most photographic lenses are made of glass and will therefore filter out most ultraviolet light. Instead, expensive lenses made of quartz must be used. Infrared films may be shot in standard cameras using an infrared band- or long-pass filters, although the infrared focal point must be compensated for.
Exposure and focusing are difficult when using UV or IR film with a camera and lens designed for visible light. The ISO standard for film speed only applies to visible light, so visual-spectrum light meters are nearly useless. Film manufacturers can supply suggested equivalent film speeds under different conditions, and recommend heavy bracketing (e.g., "with a certain filter, assume ISO 25 under daylight and ISO 64 under tungsten lighting"). This allows a light meter to be used to estimate an exposure. The focal point for IR is slightly farther away from the camera than visible light, and UV slightly closer; this must be compensated for when focusing. Apochromatic lenses are sometimes recommended due to their improved focusing across the spectrum.
Film optimized for detecting X-ray radiation is commonly used for medical radiography and industrial radiography by placing the subject between the film and a source of X-rays or gamma rays, without a lens, as if a translucent object were imaged by being placed between a light source and standard film. Unlike other types of film, X-ray film has a sensitive emulsion on both sides of the carrier material. This reduces the X-ray exposure for an acceptable imagea desirable feature in medical radiography. The film is usually placed in close contact with phosphor screen(s) and/or thin lead-foil screen(s), the combination having a higher sensitivity to X-rays. Because film is sensitive to x-rays, its contents may be wiped by airport baggage scanners if the film has a speed higher than 800 ISO. This property is exploited in Film badge dosimeters.
Film optimized for detecting X-rays and gamma rays is sometimes used for radiation dosimetry.
Film has a number of disadvantages as a scientific detector: it is difficult to calibrate for photometry, it is not re-usable, it requires careful handling (including temperature and humidity control) for best calibration, and the film must physically be returned to the laboratory and processed. Against this, photographic film can be made with a higher spatial resolution than any other type of imaging detector, and, because of its logarithmic response to light, has a wider dynamic range than most digital detectors. For example, Agfa 10E56 holographic film has a resolution of over 4,000 lines/mmequivalent to a pixel size of 0.125 micrometersand an active dynamic range of over five orders of magnitude in brightness, compared to typical scientific CCDs that might have pixels of about 10 micrometers and a dynamic range of 3–4 orders of magnitude.
Special films are used for the long exposures required by astrophotography.
Lith films used in the printing industry. In particular when exposed via a ruled-glass screen or contact-screen, halftone images suitable for printing could be generated.
Encoding of metadata
Some film cameras have the ability to read metadata from the film canister or encode metadata on film negatives.
Negative imprinting
Negative imprinting is a feature of some film cameras, in which the date, shutter speed and aperture setting are recorded on the negative directly as the film is exposed. The first known version of this process was patented in the United States in 1975, using half-silvered mirrors to direct the readout of a digital clock and mix it with the light rays coming through the main camera lens. Modern SLR cameras use an imprinter fixed to the back of the camera on the film backing plate. It uses a small LED display for illumination and optics to focus the light onto a specific part of the film. The LED display is exposed on the negative at the same time the picture is taken. Digital cameras can often encode all the information in the image file itself. The Exif format is the most commonly used format.
DX codes
In the 1980s, Kodak developed DX Encoding (from Digital indeX), or DX coding, a feature that was eventually adapted by all camera and film manufacturers. DX encoding provides information on both the film cassette and on the film regarding the type of film, number of exposures, speed (ISO/ASA rating) of the film. It consists of three types of identification. First is a barcode near the film opening of the cassette, identifying the manufacturer, film type and processing method (see image below left). This is used by photofinishing equipment during film processing. The second part is a barcode on the edge of the film (see image below right), used also during processing, which indicates the image film type, manufacturer, frame number and synchronizes the position of the frame. The third part of DX coding, known as the DX Camera Auto Sensing (CAS) code, consists of a series of 12 metal contacts on the film cassette, which beginning with cameras manufactured after 1985 could detect the type of film, number of exposures and ISO of the film, and use that information to automatically adjust the camera settings for the speed of the film.
Common sizes of film
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History
The earliest practical photographic process was the daguerreotype; it was introduced in 1839 and did not use film. The light-sensitive chemicals were formed on the surface of a silver-plated copper sheet. The calotype process produced paper negatives. Beginning in the 1850s, thin glass plates coated with photographic emulsion became the standard material for use in the camera. Although fragile and relatively heavy, the glass used for photographic plates was of better optical quality than early transparent plastics and was, at first, less expensive. Glass plates continued to be used long after the introduction of film, and were used for astrophotography and electron micrography until the early 2000s, when they were supplanted by digital recording methods. Ilford continues to manufacture glass plates for special scientific applications.
The first flexible photographic roll film was sold by George Eastman in 1885, but this original "film" was actually a coating on a paper base. As part of the processing, the image-bearing layer was stripped from the paper and attached to a sheet of hardened clear gelatin. The first transparent plastic roll film followed in 1889. It was made from highly flammable cellulose nitrate film.
Although cellulose acetate or "safety film" had been introduced by Kodak in 1908, at first it found only a few special applications as an alternative to the hazardous nitrate film, which had the advantages of being considerably tougher, slightly more transparent, and cheaper. The changeover was completed for X-ray films in 1933, but although safety film was always used for 16 mm and 8 mm home movies, nitrate film remained standard for theatrical 35 mm films until it was finally discontinued in 1951.
Hurter and Driffield began pioneering work on the light sensitivity of photographic emulsions in 1876. Their work enabled the first quantitative measure of film speed to be devised. They developed H&D curves, which are specific for each film and paper. These curves plot the photographic density against the log of the exposure, to determine sensitivity or speed of the emulsion and enabling correct exposure.
Spectral sensitivity
Early photographic plates and films were usefully sensitive only to blue, violet and ultraviolet light. As a result, the relative tonal values in a scene registered roughly as they would appear if viewed through a piece of deep blue glass. Blue skies with interesting cloud formations photographed as a white blank. Any detail visible in masses of green foliage was due mainly to the colorless surface gloss. Bright yellows and reds appeared nearly black. Most skin tones came out unnaturally dark, and uneven or freckled complexions were exaggerated. Photographers sometimes compensated by adding in skies from separate negatives that had been exposed and processed to optimize the visibility of the clouds, by manually retouching their negatives to adjust problematic tonal values, and by heavily powdering the faces of their portrait sitters.
In 1873, Hermann Wilhelm Vogel discovered that the spectral sensitivity could be extended to green and yellow light by adding very small quantities of certain dyes to the emulsion. The instability of early sensitizing dyes and their tendency to rapidly cause fogging initially confined their use to the laboratory, but in 1883 the first commercially dye-sensitized plates appeared on the market. These early products, described as isochromatic or orthochromatic depending on the manufacturer, made possible a more accurate rendering of colored subject matter into a black-and-white image. Because they were still disproportionately sensitive to blue, the use of a yellow filter and a consequently longer exposure time were required to take full advantage of their extended sensitivity.
In 1894, the Lumière Brothers introduced their Lumière Panchromatic plate, which was made sensitive, although very unequally, to all colors including red. New and improved sensitizing dyes were developed, and in 1902 the much more evenly color-sensitive Perchromo panchromatic plate was being sold by the German manufacturer Perutz. The commercial availability of highly panchromatic black-and-white emulsions also accelerated the progress of practical color photography, which requires good sensitivity to all the colors of the spectrum for the red, green and blue channels of color information to all be captured with reasonable exposure times.
However, all of these were glass-based plate products. Panchromatic emulsions on a film base were not commercially available until the 1910s and did not come into general use until much later. Many photographers who did their own darkroom work preferred to go without the seeming luxury of sensitivity to reda rare color in nature and uncommon even in human-made objectsrather than be forced to abandon the traditional red darkroom safelight and process their exposed film in complete darkness. Kodak's popular Verichrome black-and-white snapshot film, introduced in 1931, remained a red-insensitive orthochromatic product until 1956, when it was replaced by Verichrome Pan. Amateur darkroom enthusiasts then had to handle the undeveloped film by the sense of touch alone.
Introduction to color
Experiments with color photography began almost as early as photography itself, but the three-color principle underlying all practical processes was not set forth until 1855, not demonstrated until 1861, and not generally accepted as "real" color photography until it had become an undeniable commercial reality in the early 20th century. Although color photographs of good quality were being made by the 1890s, they required special equipment, separate and long exposures through three color filters, complex printing or display procedures, and highly specialized skills, so they were then exceedingly rare.
The first practical and commercially successful color "film" was the Lumière Autochrome, a glass plate product introduced in 1907. It was expensive and not sensitive enough for hand-held "snapshot" use. Film-based versions were introduced in the early 1930s and the sensitivity was later improved. These were "mosaic screen" additive color products, which used a simple layer of black-and-white emulsion in combination with a layer of microscopically small color filter elements. The resulting transparencies or "slides" were very dark because the color filter mosaic layer absorbed most of the light passing through. The last films of this type were discontinued in the 1950s, but Polachrome "instant" slide film, introduced in 1983, temporarily revived the technology.
"Color film" in the modern sense of a subtractive color product with a multi-layered emulsion was born with the introduction of Kodachrome for home movies in 1935 and as lengths of 35 mm film for still cameras in 1936; however, it required a complex development process, with multiple dyeing steps as each color layer was processed separately. 1936 also saw the launch of Agfa Color Neu, the first subtractive three-color reversal film for movie and still camera use to incorporate color dye couplers, which could be processed at the same time by a single color developer. The film had some 278 patents. The incorporation of color couplers formed the basis of subsequent color film design, with the Agfa process initially adopted by Ferrania, Fuji and Konica and lasting until the late 70s/early 1980s in the West and 1990s in Eastern Europe. The process used dye-forming chemicals that terminated with sulfonic acid groups and had to be coated one layer at a time. It was a further innovation by Kodak, using dye-forming chemicals which terminated in 'fatty' tails which permitted multiple layers to coated at the same time in a single pass, reducing production time and cost that later became universally adopted along with the Kodak C-41 process.
Despite greater availability of color film after WWII during the next several decades, it remained much more expensive than black-and-white and required much more light, factors which combined with the greater cost of processing and printing delayed its widespread adoption. Decreasing cost, increasing sensitivity and standardized processing gradually overcame these impediments. By the 1970s, color film predominated in the consumer market, while the use of black-and-white film was increasingly confined to photojournalism and fine art photography.
Effect on lens and equipment design
Photographic lenses and equipment are designed around the film to be used. Although the earliest photographic materials were sensitive only to the blue-violet end of the spectrum, partially color-corrected achromatic lenses were normally used, so that when the photographer brought the visually brightest yellow rays to a sharp focus, the visually dimmest but photographically most active violet rays would be correctly focused, too. The introduction of orthochromatic emulsions required the whole range of colors from yellow to blue to be brought to an adequate focus. Most plates and films described as orthochromatic or isochromatic were practically insensitive to red, so the correct focus of red light was unimportant; a red window could be used to view the frame numbers on the paper backing of roll film, as any red light which leaked around the backing would not fog the film; and red lighting could be used in darkrooms. With the introduction of panchromatic film, the whole visible spectrum needed to be brought to an acceptably sharp focus. In all cases a color cast in the lens glass or faint colored reflections in the image were of no consequence as they would merely change the contrast a little. This was no longer acceptable when using color film. More highly corrected lenses for newer emulsions could be used with older emulsion types, but the converse was not true.
The progression of lens design for later emulsions is of practical importance when considering the use of old lenses, still often used on large-format equipment; a lens designed for orthochromatic film may have visible defects with a color emulsion; a lens for panchromatic film will be better but not as good as later designs.
The filters used were different for the different film types.
Decline
Film remained the dominant form of photography until the early 21st century, when advances in digital photography drew consumers to digital formats. The first consumer electronic camera, the Sony Mavica was released in 1981, the first digital camera, the Fuji DS-X released in 1989, coupled with advances in software such as Adobe Photoshop which was released in 1989, improvements in consumer level digital color printers and increasingly widespread computers in households during the late 20th century facilitated uptake of digital photography by consumers.
The initial take up of digital cameras in the 1990s was slow due to their high cost and relatively low resolution of the images (compared to 35mm film), but began to make inroads among consumers in the point and shoot market and in professional applications such as sports photography where speed of results including the ability to upload pictures direct from stadia was more critical for newspaper deadlines than resolution. A key difference compared to film was that early digital cameras were soon obsolete, forcing users into a frequent cycle of replacement until the technology began to mature, whereas previously people might have only owned one or two film cameras in their lifetime. Consequently, photographers demanding higher quality in sectors such as weddings, portraiture and fashion where medium format film predominated were the last to switch once resolution began to reach acceptable levels with the advent of 'full frame' sensors, 'digital backs' and medium format digital cameras.
Film camera sales based on CIPA figures peaked in 1998, before declining rapidly after 2000 to reach almost zero by the end of 2005 as consumers switched en masse to digital cameras (sales of which subsequently peaked in 2010). These changes foretold a similar reduction in film sales. Figures for Fujifilm show global film sales, having grown 30% in the preceding five years, peaked around the year 2000. Film sales then began a period of year-on-year falling sales, of increasing magnitude from 2003 to 2008, reaching 30% per annum before slowing. By 2011, sales were less than 10% of the peak volumes. Similar patterns were experienced by other manufacturers, varying by market exposure, with global film sales estimated at 200 million rolls in 1999 declining to only 5 million rolls by 2009. This period wreaked havoc on the film manufacturing industry and its supply chain optimised for high production volumes, plummeting sales saw firms fighting for survival. Agfa-Gevaert's decision to sell off its consumer facing arm (Agfaphoto) in 2004, was followed by a series of bankruptcies of established film manufacturers: Ilford Imaging UK in 2004, Agfaphoto in 2005, Forte in 2007, Foton in 2007, Polaroid in 2001 and 2008, Ferrania in 2009, and Eastman Kodak in 2012 (the latter only surviving after massive downsizing whilst Ilford was rescued by a management buyout). Konica-Minolta closed its film manufacturing business and exited the photographic market entirely in 2006, selling its camera patents to Sony, and Fujifilm successfully moved to rapidly diversify into other markets. The impact of this paradigm shift in technology subsequently rippled though the downstream photo processing and finishing businesses.
Although modern photography is dominated by digital users, film continues to be used by enthusiasts. Film remains the preference of some photographers because of its distinctive "look".
Renewed interest in recent years
Despite the fact that digital cameras are by far the most commonly-used photographic tool and that the selection of available photographic films is much smaller than it once was, sales of photographic film have been on a steady upward trend. Kodak (which was under bankruptcy protection from January 2012 to September 2013) and other companies have noticed this upward trend: Dennis Olbrich, President of the Imaging Paper, Photo Chemicals and Film division at Kodak Alaris, has stated that sales of their photographic films have been growing over the past three or four years. UK-based Ilford have confirmed this trend and conducted extensive research on this subject matter, their research showing that 60% of current film users had only started using film in the past five years and that 30% of current film users were under 35 years old. Annual film sales, which were estimated to reach a low of 5 million rolls in 2009, have since doubled to around 10 million rolls in 2019. A key challenge for the industry is that production relies on the remaining coating facilities that were built for the peak years of demand, but as demand has grown capacity constraints in some of the other process steps which have been downscaled, such as converting film, have caused production bottlenecks for companies such as Kodak.
In 2013 Ferrania, an Italy-based film manufacturer which ceased production of photographic films between the years 2009 and 2010, was acquired by the new Film Ferrania S.R.L taking over a small part of the old company's manufacturing facilities using its former research facility, and re-employed some workers who had been laid off three years earlier when the company stopped production of film.
In November of the same year, the company started a crowdfunding campaign with the goal of raising $250,000 to buy tooling and machines from the old factory, with the intention of putting some of the films that had been discontinued back into production, the campaign succeeded and in October 2014 was ended with over $320,000 being raised. In February 2017, Film Ferrania unveiled their "P30" 80 ASA, Panchromatic black and white film, in 35mm format.
Kodak announced on January 5, 2017, that Ektachrome, one of Kodak's most well known transparency films that had been discontinued between 2012 and 2013, would be reformulated and manufactured once again, in 35 mm still and Super 8 motion picture film formats. Following the success of the release, Kodak expanded Ektachrome's format availability by also releasing the film in 120 and 4x5 formats.
Japan-based Fujifilm's instant film "Instax" cameras and paper have also proven to be very successful, and have replaced traditional photographic films as Fujifilm's main film products, while they continue to offer traditional photographic films in various formats and types.
Reusable film
In 2023, a Finnish chemist Sami Vuori invented a reusable film that uses synthetic hackmanite (Na8Al6Si6O24(Cl,S)2) as the photosensitive medium. The film contains small hackmanite particles that color purple with exposure to ultraviolet radiation (e.g. 254 nm), after which the film is loaded into the camera. Visible light bleaches the hackmanite particles back to white, which gives rise to the formation of a positive image. The film can then be scanned with a typical document scanner and then colored again with UV. If the user wants to spare the image, the film can be put into a dark place, as the bleaching process stops completely in the absence of light.
On top of reusability and not needing any developing or chemicals, another advantage with this type of photochromic film is that it does not need gelatin, which renders it a vegan alternative. However, the main disadvantage is still the film's very slow exposure, requiring hours of exposure time. This means that currently this type of film can be used only in ultra-long-exposure film photography where the subject is e.g. a city center where the photographer wants to fade all movement.
Another reusable film invented by Liou et al. is based on 9-methylacridinium-intercalated clay particles, but erasing the image requires dipping the material in sulfuric acid.
Image gallery
| Technology | Optical instruments | null |
21561194 | https://en.wikipedia.org/wiki/Fall%20armyworm | Fall armyworm | The fall armyworm (Spodoptera frugiperda) is a species in the order Lepidoptera and one of the species of the fall armyworm moths distinguished by their larval life stage. The term "armyworm" can refer to several species, often describing the large-scale invasive behavior of the species' larval stage. It is regarded as a pest and can damage and destroy a wide variety of crops, which causes large economic damage. Its scientific name derives from frugiperda, which is Latin for lost fruit, named because of the species' ability to destroy crops. Because of its propensity for destruction, the fall armyworm's habits and possibilities for crop protection have been studied in depth. It is also a notable case for studying sympatric speciation, as it appears to be diverging into two species currently. Another remarkable trait of the larva is that they consistently practice cannibalism, despite its fitness costs.
The fall armyworm is active at a different time of year from the true armyworm, another species in the order Lepidoptera and family Noctuidae, but of the genus Mythimna. Outbreaks of the true armyworm usually occur during the early part of the summer; the fall armyworm does most damage in the late summer in the southern part of the United States, and early fall in the northern regions.
Description
The adult moths are wing tip to wing tip, with a brown or gray forewing, and a white hindwing. There is slight sexual dimorphism, with males having more patterns and a distinct white spot on each of their forewings. The first larval instar is light colored with a larger dark head. As they develop through instars, they become browner with white lengthwise lines. They also develop dark spots with spines.
Geographic range
Native range
The fall armyworm is widely distributed in eastern and central North America and in South America. It cannot survive overwinter in below freezing temperatures, so it only survives the winter in the most southern regions of the United States, namely Texas and Florida. Because of this, the fall armyworm is a more prominent pest in southeastern states. However, seasonally it will spread across the eastern United States and up to southern Canada, inhabiting areas with suitable food supplies.
Introduced range
The potential global distribution of S. frugiperda has been modelled using CLIMEX. The modelled global potential distribution reflects the marked seasonal range dynamics experienced in North America, with much of the potential range in Europe, South Africa, China and Australia consisting of habitat that is only climatically suitable during the warmer months.. A more recent physiologically-based population dynamics model was developed for assessing the potential distribution of S. frugiperda in Europe. The model showed that the Mediterranean coastal areas of Southern Europe might be particularly suitable for the establishment of the species.
Africa
S. frugiperda was first found on the African continent in 2013 in Sao Tome, then spread through Nigeria, Benin, Togo, and was found in Ghana in February 2017. In December 2020 S. frugiperda was first found in Syria in Daraa on the Jordanian border and is believed to have arrived from there without human assistance, having just been found in that country also. The fall armyworm is causing significant damage to maize crops in Africa and has great potential for further spread and economic damage. It has since spread to 28 countries in Africa.
Asia
S. frugiperda was first detected in Bangladesh in late 2018. it has reached 37 districts. As a result of the introduction of S. frugiperda and Lumpy Skin Disease within a few months of each other, the FAO, the World Food Programme, Bangladesh Government officials, and others agreed to begin improving Bangladesh's agricultural emergency response capabilities. The use of two biopesticides – Spodoptera frugiperda nuclear polyhedrosis virus/SfNPV (the SNPV/single nuclear polyhedrosis virus specific to S. frugiperda) and Habrobracon hebetor – is recommended.
In December 2018, the Fall armyworm began to spread widely in India. In January 2019, a heavy infestation of fall armyworm was recorded in corn plantations in Sri Lanka.
The pest was first detected in China in the southwest province of Yunnan in January 2019 (or June 2019). Through 2019, the pest infested a total of 26 provinces. The armyworm is expected in 2020 to hit China's Northeast wheat belt. A report issued by the Ministry of Agriculture and Rural Affairs rates the situation as "very grave".
The fall armyworm was first reported in Southeast Asia in late 2018 in Thailand and Myanmar and its presence is now confirmed in almost all Southeast Asian countries.
Oceania
In January 2020 S. frugiperda was detected on the Torres Strait Islands, in February in North Queensland, and then continued into the rest of Queensland, and the Northern Territory, Western Australia, and then in September was found in New South Wales between Moree and Boggabilla (and later in Narrabri, Wee Waa, Dubbo, Breeza, and Maitland). S. frugiperda is expected to severely impact Queensland's wool industry because it feeds on all major grazing plants. It was observed in traps baited with a male pheromone lure, firstly on Darnley Island and Saibai islands in the Torres Strait, and subsequently on the mainland near Croydon. Within a week it was officially declared ineradicable. In April 2020, it was detected in Papua New Guinea, spreading across the Torres Strait.
Fall armyworm was first detected in New Zealand in February 2022. Biosecurity New Zealand and sector partners ran a biosecurity response to limit the spread of Fall armyworm and try to eradicate it from New Zealand. This included surveillance and research to better understand the moth, its spread, and potential impacts in New Zealand. By April 2023, it became clear that Fall armyworm was widespread, particularly in the North Island, and that eradication was unlikely because it had been windblown from Australia, and this is likely to repeatedly occur.
Food resources
Caterpillars
The armyworm's diet consists mainly of grasses and grain crops such as corn, but the species has been noted to consume over 80 different plants (50 non-economic and 30 economic plants). Armyworms earned their common name by eating all plant matter they encounter in their wide dispersals, like a large army. A few sweet corn varieties have partial, but not complete, resistance to armyworms. The resistance comes from a unique 33-kD proteinase that the corn produces when it is being fed on by fall armyworms or other larvae. This protein was found to significantly decrease fall armyworm larva growth.
Cannibalism
When possible, larvae will cannibalize the larvae of smaller instars. A 1999 study showed that cannibalism only benefits the caterpillar when other food is scarce. Despite this, the caterpillars will cannibalize others whenever they can, even though it was found to decrease their own fitness in many cases. One known reason why cannibalism is detrimental to the fall armyworm is because of disease transmission to the cannibal. In nature, the negative effects of cannibalism may be balanced by the fact that cannibalism removes competitors, thereby making more resources accessible and indirectly increasing fall armyworms' fitness.
Adults
Adult moths sip nectar from flowers such as that of witch hazel (Hamamelis virginiana).
Life history
The fall armyworm's life cycle is completed within 30 days during summer, and 60 days during the spring and autumn seasons; during the winter, these caterpillars' life cycle lasts about 80 to 90 days.[3] The number of generations a moth will have in a year varies based on climate, but in her life span a female will typically lay about 1,500 eggs. Because larvae cannot enter into diapause they cannot survive cold temperatures.
Egg
The armyworm's egg is dome-shaped, and measures around in diameter and in height. Females prefer to lay eggs on the underside of leaves, but in high populations they will lay them just about anywhere. In warm weather, the eggs will hatch into larvae within a few days.
Larva
The larvae go through six different instars, each varying slightly in physical appearance and pattern. The larva process lasts from 14 to 30 days, again depending on temperatures. The mature caterpillar is about in length. This is the most destructive life stage as the larvae have biting mouth parts. The larvae have a distinctive inverted Y suture on the forehead.
Pupa
The larvae then pupate underground for 7 to 37 days in a cocoon they form of soil and silk. Duration and survival of the pupal stage depend on the temperature of the environment.
Adults
Once emerged, the adults live for about 10 days, and sometimes up to 21 days, with the female laying most of her eggs early in life. Adults are nocturnal and fare best during warm and humid nights.
Migration
Adults are capable of flying long distances, so even though they are unable to overwinter north of the southern region of the United States, the moths can migrate as far north as southern Canada in warm months. Their migration rate is remarkably fast, estimated at per generation. Some scientists speculate that this fast migration is aided by the movement of air in weather fronts.
Neurochemistry
Allatotropin and allatotropin+allatostatin C – neuropeptides – extracted from Manduca sexta were both found to suppress feeding in all life stages, increase larval mortality, and reduce adult lifespan, by Oeh et al 2000.
Enemies
Predators
Fall armyworm caterpillars are directly preyed upon by many invertebrates and vertebrates. Common predators include birds, rodents, beetles, earwigs, and other insects. It has been shown that direct predation can cause significant losses to caterpillar populations. The larva's main defense against enemies is their ability to reach large numbers and migrate before seasonal conditions are suitable for predators.
Parasitoids
Fly and wasp parasitoids target the fall armyworm, most commonly Archytas marmoratus, Cotesia marginiventris, and Chelonus texanus. The armyworm is also vulnerable to additional parasitoids, varying with location. In 2018, egg parasitoid wasps of the genera Telenomus and Trichogramma were discovered to attack army worm eggs in East Africa.
Cotesia icipe is another African braconid wasp suitable for the biological control of this lepidoptera.
Parasites and disease
Fifty-three different parasite species have been discovered in fall armyworm larvae, spanning ten different families. Often larvae can survive through much of their crop consumption despite outbreaks of disease, because of the larva's fast life cycle. Despite this, parasites of the fall armyworm are being studied extensively as a means of fighting armyworm attacks on crops. One suggested approach would be to introduce parasites from South America to North American fall armyworms, and vice versa.
Fungi
In February 2021, it was reported that an Australian agronomist Georgia Rodger had found at a property near Beaudesert (southern Queensland) the tropical fungus Nomuraea rileyi which was known to be effective in killing and consuming fall armyworms. Samples of this were sent to Maree Crawford, the insect pathologist at the Queensland Department of Agriculture for further analysis. Australian entomologists have said the finding is reassuring and that laboratory tests have been promising. This is substantiated by various studies including a 2018 journal article which looked into the effectiveness of N. rileyi had on infestations of armyworms in Indian maize crops. The study concluded N. rileyi could potentially be a cost-effective tool in combating the pest, compatible with eco-friendly management practices, although further studies were required. Farmers in Australia have struggled to control the pest which has been destroying crops, prompting concerns about potential food shortages which could cause an increase in food prices for consumers. The N. rileyi research has given them hope that this can be avoided.
Subspecies
The fall armyworm may be presently undergoing a divergence into two separate species. These two strains have major genetic differences that are connected to the plants they feed on, even though both still exist in the same area (sympatric speciation). These two strains can be loosely categorized into a rice strain and a corn strain. This separation is occurring because of differences in habitat (preferred host plant), and differences in reproductive behavior. The reproductive differences can be divided into two categories: difference in the timing of mating at night, and difference in female sex pheromones.
Mating
Mate searching behavior and male–male conflict
A female attracts males by perching atop the host plant feeding area and releasing a sex pheromone as the signal that she wishes to mate. The pheromone has been studied and found to contain the components Z7-12 and Z9-14. Each female only mates once per night; this creates a physical conflict between the multiple males that will fly towards a ready female. There is an order to which the females call and mate: virgin females do first, females who have mated once next, and females who have already mated multiple times call and mate last during the night.
Interactions with humans
Research use
S. frugiperda cells (Sf9 and Sf21 cell lines) are commonly used in biomedical research for the purpose of recombinant protein expression using insect-specific viruses called baculoviruses.
Pest of crop plants
Because of their food preferences, fall armyworm larvae can wreak havoc on a wide range of crops. The first historical account of the fall armyworm's destruction was in 1797 in Georgia. Destruction can happen almost over night, because the first stages of a caterpillar's life require very little food, and the later stages require about 50 times more. Because of this rapid change in food consumption, the presence of larvae will not be noticed until they have destroyed almost everything in as little as a night. Some examples of targeted crops include cotton, tobacco, sweet corn, rice, peanuts, and even fruits such as apples, oranges, and many more. The list of possible food sources for the worms is extensive, so crop damage is wide-ranging. It is estimated that almost 40 percent of those species that armyworms target are economically important. Because the larvae eat so much of the plant, they are very detrimental to crop survival and yield. In corn, larvae will even burrow into the corn ear to eat the kernels.
The UN Food and Agriculture Organization estimates that S. frugiperda will reduce maize/corn yields by /annum if not successfully controlled. The fall armyworm have proved to be a pest in many regions, and methods of control continue to be developed.
Africa
The fall armyworm was identified in Africa in 2016. In early 2017, armyworms infested large swathes of corn crops across southern Africa, devastating the livelihoods of many farmers. It is thought they arrived as an invasive species from the Americas as eggs in imported produce. This is causing immense concern among agricultural experts, due to the potentially huge amount of damage this invasive species will do to African food crops if allowed to spread. Many African countries have agreed to take urgent actions against armyworms.
Sri Lanka
After being first reported in India in May 2018 in Tamil Nadu, then the Sri Lankan Ministry of Agriculture issued a warning notice to farmers in the northwestern and north central provinces about possible fall armyworm invasion. At the time of warning, crop destruction had already been reported from the Ampara, Anuradhapura, and Polonnaruwa areas. The larvae are known among the local people as Sena dalambuwa (armyworm caterpillar). Not only corn, but also sugarcane plantations were attacked by the caterpillars in Anuradhapura, Ampara, and Monaragala districts.
In December 2018, heavy infestations in corn cultivation were identified. The spread of the moth leads to attack corn all around the country within weeks. On 6 January 2019, caterpillars spread to the Monaragala district and devastated corn crops. At the end of January 2019, the armyworm was present in all districts of Sri Lanka except Nuwara Eliya and Jaffna.
On 29 December 2018, armyworms were recorded from paddy cultivations in the Sinhapura area of Polonnaruwa. In January 2019, caterpillars were also recorded from paddy cultivations of the Nochchiyagama area in the Anuradhapura district.
The Sri Lankan Department of Agriculture recommended 12 pesticides under three categories, to be used alternately every seven days. Organic farming expert, Thilak Kandegama said that the threat can be overcome by sprinkling rice husk ashes as a repellent. Agricultural Ministry also decided to use drone technology for the spraying of insecticides to control the spreading of caterpillars.
Management and control
Because of the fall armyworms' great destructive power, farmers must go to great lengths to deter the larvae. Insecticide is a widely used form of protection; in southern regions, farmers may have to apply insecticide to corn every day. Agricultural drones have been used to apply pesticides, used in China, Vietnam, Zambia and other regions.
The CABI-led programme, Plantwise and partners have several recommendations for managing fall armyworm, these include: planting early, avoiding staggered planting, and inter-cropping with crops that are not susceptible to fall armyworm, such as cassava or yam. They also recommend conserving shelters and flowering plants on the edges for beneficial insects such as ground beetles and parasitoids.
Inter-cropping with the "push-pull" technique with crops such as Desmodium and Napier grass can be used to control fall armyworm.
For some crops, including wheat, sorghum, millet and rice, it is recommend by Plantwise partners to plant short maturing and varieties that are less preferred by S. frugiperda.
Another strategy is to plant crops earlier to avoid the increase in armyworm numbers as the summer progresses.
In South Africa, farmers are using pheromone lures with a combination of Dichlorvos blocks to trap and eliminate male armyworms, with the intention of disrupting mating cycles.
CIMMYT and its partners are using forward genetics to breed for better S. frugiperda resistance in maize. Genome-wide association studies (GWAS) are the most effective method for associating S. f. resistance to the responsible genomic region, especially used in maize/corn but also wheat, sorghum, millet, rice, and legumes. The first uses of conventional breeding in the first decade of the 1900s were reported by Gernet 1917 and Hinds 1914, improving resistance in maize/corn, sorghum, millet, Cynodon dactylon, and Arachis hypogaea.
In Australia, a caterpillar-specific virus packaged as Fawligen biopesticide was approved under emergency regulations in 2020 to help control the armyworm, and the parasitoid wasp Trichogramma pretiosum is also used.
Directorate of plant protection Quarantine and storage, Ministry of Agriculture, Govt of India regularly issues advisories from time to time to manage the menace of Fall Army Worm in India.
| Biology and health sciences | Lepidoptera | Animals |
3907049 | https://en.wikipedia.org/wiki/Lesser%20florican | Lesser florican | The lesser florican (Sypheotides indicus), also known as the likh or kharmore, is the smallest in the bustard family and the only member of the genus Sypheotides. It is endemic to the Indian Subcontinent where it is found in tall grasslands and is best known for the leaping breeding displays made by the males during the monsoon season. The male has a contrasting black and white breeding plumage and distinctive elongated head feathers that extend behind the neck. These bustards are found mainly in northwestern and central India during the summer but are found more widely distributed across India in winter. The species is highly endangered and has been extirpated in some parts of its range such as Pakistan. It is threatened both by hunting and habitat degradation. The only similar species is the Bengal florican (Houbarobsis bengalensis) which is larger and lacks the white throat, collar and elongated plumes.
Taxonomy and systematics
In 1782 the English illustrator John Frederick Miller included a hand-coloured plate of a female lesser florican in his Icones animalium et plantarum. He coined the binomial name Otis indica. It is now the only species placed in the genus Sypheotides that was introduced in 1839 by the French naturalist René Lesson. The species is monotypic: no subspecies are recognised.
The two species of smaller bustards have been called "floricans". The word has been thought to be of Dutch origin. The genus Sypheotides earlier included what is now Houbaropsis bengalensis (or Bengal florican), the two species being small and showing reverse sexual size dimorphism. The tarsus is long in Sypheotides and the seasonal plumage change in male has led to the retention of the separate genus, although the two genera are evolutionarily close. Male and female plumages were initially thought of as separate species leading to the names aurita and indica and the species has been placed in the past in the genera Otis, Eupodotis and Sypheotis. The species ending which is related to the gender of the Latin genus has been debated and it believed that indicus is correct.
The horizontal body carriage, size and habit of holding up their tail feathers when walking on the ground have led their local names to make associations with peacocks, with a popular name being the equivalent of "grass peacock" (such as khar-mor, tan-mor) in some areas. the name Likh is used in northwestern India and adopted by British sportsmen in India.
Description
A male in breeding plumage has a black head, neck and lower parts. However, his throat is white. Around three 4 inch long, ribbon-like feathers arise from behind the ear-coverts on each side of the head and extend backwards, curving up and ending in spatulate tip. The back and scapulars are mottled in white with V-shaped marks. The wing coverts are white. After the breedings season, the male tends to have some white in the wing. The female is slightly larger than the male. The females and males in non breeding plumage are buff with black streaks with darker markings on the head and neck. The back is mottled and barred in black. The neck and upper breast are buff with the streaks decreasing towards the belly. The outer primaries of the males are thin and notched on the inner-web. The leg are pale yellow and the iris is yellow.
Young birds have a distinct U-shaped mark on the neck near the throat.
Distribution and habitat
The species was formerly more widespread across much of Indian Sub-continent, but not in Sri Lanka. It breeds mainly in the central and western parts of India. Historic records exist from the Makran coast of Balochistan province in Pakistan. A record from Burma has been questioned. The species is said to move in response to rainfall and their presence at locations can be erratic, with sudden large numbers in some seasons. About 500 males in Gujarat were ringed and nearly 18 were recovered, most of them within about 50 kilometres of their ringing sites. The preferred habitat is grasslands but it sometimes occurs in fields such as those of cotton and lentils. Breeding areas are today restricted mainly to Gujarat, Madhya Pradesh, some areas in southern Nepal and parts of Andhra Pradesh.
Managing florican habitats as grassland interspersed with croplands and pastures spared rotationally provided optimal results at low production-level.
Behaviour and ecology
These bustards are found either singly or in pairs in thick grassland or sometimes in crop fields. Indigenous tribal hunters regularly shot the males during the breeding season, as they were easy to spot because of their courtship display. It was said to be good for eating but considered inferior to the meat of the Bengal florican. They fly faster than other bustards and give a duck-like impression in flight.
Food and feeding
Lesser floricans feed on a wide variety of small vertebrates and invertebrates which include worms, centipedes, lizards, frogs and insects such as locusts, flying ants and hairy caterpillars. They are also known to feed on shoots and seeds, herbs and berries.
Usually floricans feed during the early hours of mornings or in the evenings, except in the case of newly migrated birds which feed throughout the day.
Breeding
The breeding season varies with the onset of the Southwest Monsoon and is September to October in northern India and April to May in parts of southern India.
During the breeding season, males leap suddenly from the grass with a peculiar croaking or knocking call, flutter their wings and fall back with slightly open wings. At the apogee of the leap the neck is arched backwards and the legs folded as if in a sitting posture. These jumps are repeated after intervals of about three or more minutes. The displays are made mainly in the early mornings and late evenings, but during other parts of the day in cloudy weather.
The breeding system is said to be a dispersed lek with each male holding a territory of about 1-2 hectares. Males are said to favour particular display sites and shooting of these displaying birds has led to sharp declines in the populations in the past. Lek sites tend to have flat ground with low vegetation and good visibility and well used sites usually show signs of trampling.
Females have a defensive display at nest which involves spreading their wings, tail and neck feathers. The females are said to produce a whistling call which attracts males. Males are aggressive towards other males in the neighbourhood. The nest is a shallow scrape on the ground and 3-4 (1.88 x 1.6 inches) eggs are laid. The nest location is usually in dense grass. Females take sole part in incubation and rearing the chicks. The incubation period is about 21 days.
| Biology and health sciences | Basics | Animals |
23046498 | https://en.wikipedia.org/wiki/Andalusian%20horse | Andalusian horse | The Andalusian, also known as the Pure Spanish Horse or PRE (), is a horse breed from the Iberian Peninsula, where its ancestors have lived for thousands of years. The Andalusian has been recognized as a distinct breed since the 15th century, and its conformation has changed very little over the centuries. Throughout its history, it has been known for its prowess as a war horse, and was prized by the nobility. The breed was used as a tool of diplomacy by the Spanish government, and kings across Europe rode and owned Spanish horses. During the 19th century, warfare, disease and crossbreeding reduced herd numbers dramatically, and despite some recovery in the late 19th century, the trend continued into the early 20th century. Exports of Andalusians from Spain were restricted until the 1960s, but the breed has since spread throughout the world, despite their low population. In 2010, there were more than 185,000 registered Andalusians worldwide.
Strongly built, and compact yet elegant, Andalusians have long, thick manes and tails. Their most common coat color is gray, although they can be found in many other colors. They are known for their intelligence, sensitivity and docility. A sub-strain within the breed known as the Carthusian, is considered by breeders to be the purest strain of Andalusian, though there is no genetic evidence for this claim. The strain is still considered separate from the main breed however, and is preferred by breeders because buyers pay more for horses of Carthusian bloodlines. There are several competing registries keeping records of horses designated as Andalusian or PRE, but they differ on their definition of the Andalusian and PRE, the purity of various strains of the breed, and the legalities of stud book ownership. At least one lawsuit is in progress , to determine the ownership of the Spanish PRE stud book.
The Andalusian is closely related to the Lusitano of Portugal, and has been used to develop many other breeds, especially in Europe and the Americas. Breeds with Andalusian ancestry include many of the warmbloods in Europe as well as western hemisphere breeds such as the Azteca. Over its centuries of development, the Andalusian breed has been selected for athleticism and stamina. The horses were originally used for classical dressage, driving, bullfighting, and as stock horses. Modern Andalusians are used for many equestrian activities, including dressage, show jumping and driving. The breed is also used extensively in movies, especially historical pictures and fantasy epics.
Characteristics
Andalusians stallions and geldings average at the withers and in weight; mares average and . The Spanish government has set the minimum height for registration in Spain at for males and for mares – this standard is followed by the Association of Purebred Spanish Horse Breeders of Spain (Asociación Nacional de Criadores de Caballo de Pura Raza Española or ANCCE) and the Andalusian Horse Association of Australasia. The Spanish legislation also requires that in order for animals to be approved as either "qualified" or "élite" breeding stock, stallions must stand at least and mares at least .
Andalusian horses are elegant and strongly built with a straight or slightly convex profile. Ultra convex and concave profiles are discouraged in the breed, and are penalized in breed shows. Necks are long and broad, running to well-defined withers and a massive chest. They have a short back and broad, strong hindquarters with a well-rounded croup. The breed tends to have clean legs, with no propensity for blemishes or injuries, and energetic gaits. The mane and tail are thick and long, but the legs do not have excess feathering. Andalusians tend to be docile, while remaining intelligent and sensitive. When treated with respect they are quick to learn, responsive, and cooperative.
There are two additional characteristics unique to the Carthusian strain, believed to trace back to the strain's foundation stallion Esclavo. The first is warts under the tail, a trait which Esclavo passed to his offspring, and a trait which some breeders felt was necessary to prove that a horse was a member of the Esclavo bloodline. The second characteristic is the occasional presence of "horns", which are frontal bosses, possibly inherited from Asian ancestors. The physical descriptions of the bosses vary, ranging from calcium-like deposits at the temple to small horn-like protuberances near or behind the ear. However, these "horns" are not considered proof of Esclavo descent, unlike the tail warts.
In the past, most coat colors were found, including spotted patterns. Today most Andalusians are gray or bay; in the US, around 80 percent of all Andalusians are gray. Of the remaining horses, approximately 15 percent are bay and 5 percent are black, dun or palomino or chestnut. Other colors, such as buckskin, pearl, and cremello, are rare, but are recognized as allowed colors by registries for the breed.
In the early history of the breed, certain white markings and whorls were considered to be indicators of character and good or bad luck. Horses with white socks on their feet were considered to have good or bad luck, depending on the leg or legs marked. A horse with no white markings at all was considered to be ill-tempered and vice-ridden, while certain facial markings were considered representative of honesty, loyalty and endurance. Similarly, hair whorls in various places were considered to show good or bad luck, with the most unlucky being in places where the horse could not see them – for example the temples, cheek, shoulder or heart. Two whorls near the root of the tail were considered a sign of courage and good luck.
The movement of Andalusian horses is extended, elevated, cadenced and harmonious, with a balance of roundness and forward movement. Poor elevation, irregular tempo, and excessive winging (sideways movement of the legs from the knee down) are discouraged by breed registry standards. Andalusians are known for their agility and their ability to learn difficult moves quickly, such as advanced collection and turns on the haunches. A 2001 study compared the kinematic characteristics of Andalusian, Arabian and Anglo-Arabian horses while moving at the trot. Andalusians were found to overtrack less (the degree to which the hind foot lands ahead of the front hoof print) but also exhibit greater flexing of both fore and hind joints, movement consistent with the more elevated way of going typically found in this breed. The authors of the study theorized that these characteristics of the breed's trot may contribute to their success as a riding and dressage horse.
A 2008 study found that Andalusians experience ischaemic (reduced blood flow) diseases of the small intestine at a rate significantly higher than other breeds; and stallions had higher numbers of inguinal hernias, with risk for occurrence 30 times greater than other breeds. At the same time, they also showed a lower incidence of large intestinal obstruction. In the course of the study, Andalusians also showed the highest risk of laminitis as a medical complication related to the intestinal issues.
History
Early development
The Andalusian horse is descended from the Iberian horses of Spain and Portugal, and derives its name from its place of origin, the Spanish region of Andalusia. Cave paintings show that horses have been present on the Iberian Peninsula as far back as 20,000 to 30,000 BCE. Although Portuguese historian Ruy d'Andrade hypothesized that the ancient Sorraia breed was an ancestor of the Southern Iberian breeds, including the Andalusian, genetic studies using mitochondrial DNA show that the Sorraia is part of a genetic cluster that is largely separated from most Iberian breeds.
Throughout history, the Iberian breeds have been influenced by many different peoples and cultures who occupied Spain, including the Celts, the Carthaginians, the Romans, various Germanic tribes and the Arabs. The Iberian horse was identified as a talented war horse as early as 450 BCE. Mitochondrial DNA studies of the modern Andalusian horse of the Iberian Peninsula and Barb horse of North Africa present convincing evidence that both breeds crossed the Strait of Gibraltar and were used for breeding with each other, influencing one another's bloodlines. Thus, the Andalusian may have been the first European "warmblood", a mixture of heavy European and lighter Oriental horses. Some of the earliest written pedigrees in recorded European history were kept by Carthusian monks, beginning in the 13th century. Because they could read and write, and were thus able to maintain careful records, monastics were given the responsibility for horse breeding by certain members of the nobility, particularly in Spain. Andalusian stud farms for breeding were formed in the late 15th century in Carthusian monasteries in Jerez, Seville and Cazalla.
The Carthusians bred powerful, weight-bearing horses in Andalusia for the Crown of Castile, using the finest Spanish Jennets as foundation bloodstock. These horses were a blend of Jennet and warmblood breeding, taller and more powerfully built than the original Jennet. By the 15th century, the Andalusian had become a distinct breed, and was being used to influence the development of other breeds. They were also noted for their use as cavalry horses. Even though in the 16th and 17th centuries Spanish horses had not reached the final form of the modern Andalusian, by 1667 William Cavendish, the Duke of Newcastle, called the Spanish horse of Andalusia the "princes" of the horse world, and reported that they were "unnervingly intelligent". The Iberian horse became known as the "royal horse of Europe" and was seen at many royal courts and riding academies, including those in Austria, Italy, France and Germany. By the 16th century, during the reigns of Charles V (1500–1558) and Phillip II (1556–1581), Spanish horses were considered the finest in the world. Even in Spain, quality horses were owned mainly by the wealthy. During the 16th century, inflation and an increased demand for harness and cavalry horses drove the price of horses extremely high. The always expensive Andalusian became even more so, and it was often impossible to find a member of the breed to purchase at any price.
Dissemination
Spanish horses also were spread widely as a tool of diplomacy by the government of Spain, which granted both horses and export rights to favored citizens and to other royalty. As early as the 15th century, the Spanish horse was widely distributed throughout the Mediterranean, and was known in northern European countries, despite being less common and more expensive there. As time went on, kings from across Europe, including every French monarch from Francis I to Louis XVI, had equestrian portraits created showing themselves riding Spanish-type horses. The kings of France, including Louis XIII and Louis XIV, especially preferred the Spanish horse; the head groom to Henri IV, Salomon de la Broue, said in 1600, "Comparing the best horses, I give the Spanish horse first place for its perfection, because it is the most beautiful, noble, graceful and courageous". War horses from Spain and Portugal began to be introduced to England in the 12th century, and importation continued through the 15th century. In the 16th century, Henry VIII received gifts of Spanish horses from Charles V, Ferdinand II of Aragon and the Duke of Savoy and others when he wed Katherine of Aragon. He also purchased additional war and riding horses through agents in Spain. By 1576, Spanish horses made up one third of British royal studs at Malmesbury and Tutbury. The Spanish horse peaked in popularity in Great Britain during the 17th century, when horses were freely imported from Spain and exchanged as gifts between royal families. With the introduction of the Thoroughbred, interest in the Spanish horse faded after the mid-18th century, although they remained popular through the early 19th century. The Conquistadors of the 16th century rode Spanish horses, particularly animals from Andalusia, and the modern Andalusian descended from similar bloodstock. By 1500, Spanish horses were established in studs on Santo Domingo, and Spanish horses made their way into the ancestry of many breeds founded in North and South America. Many Spanish explorers from the 16th century on brought Spanish horses with them for use as war horses and later as breeding stock. By 1642, the Spanish horse had spread to Moldavia, to the stables of Transylvanian prince George Rakoczi.
19th century to present
Despite their ancient history, all living Andalusians trace to a small number of horses bred by religious orders in the 18th and 19th centuries. An influx of heavy horse blood beginning in the 16th century, resulted in the dilution of many of the bloodlines; only those protected by selective breeding remained intact to become the modern Andalusian. During the 19th century, the Andalusian breed was threatened because many horses were stolen or requisitioned in wartime, including the War of the Oranges, the Peninsular War and the three Carlist Wars. Napoleon's invading army also stole many horses. One herd of Andalusians was hidden from the invaders however, and subsequently used to renew the breed. In 1822, breeders began to add Norman blood into Spanish bloodlines, as well as further infusions of Arabian blood. This was partially because increasing mechanization and changing needs within the military called for horses with more speed in cavalry charges as well as horses with more bulk for pulling gun carriages. In 1832, an epidemic seriously affected Spain's horse population, from which only one small herd survived in a stud at the monastery in Cartuja. During the 19th and early 20th centuries, European breeders, especially the Germans, changed from an emphasis on Andalusian and Neapolitan horses (an emphasis that had been in place since the decline of chivalry), to an emphasis on the breeding of Thoroughbreds and warmbloods, further depleting the stock of Andalusians. Despite this change in focus, Andalusian breeding slowly recovered, and in 1869, the Seville Horse Fair (originally begun by the Romans), played host to between ten and twelve thousand Spanish horses. In the early 20th century, Spanish horse breeding began to focus on other breeds, particularly draft breeds, Arabians, Thoroughbreds and crosses between these breeds, as well as crosses between these breeds and the Andalusian. The purebred Andalusian was not viewed favorably by breeders or the military, and their numbers decreased significantly.
Andalusians only began to be exported from Spain in 1962. The first Andalusians were imported into Australia in 1971, and in 1973 the Andalusian Horse Association of Australasia was formed for the registration of these Andalusians and their offspring. Strict quarantine guidelines prohibited the importation of new Andalusian blood to Australia for many years, but since 1999, regulations have been relaxed and more than half a dozen new horses have been imported. Bloodines in the United States also rely on imported stock, and all American Andalusians can be traced directly to the stud books in Portugal and Spain. There are around 8,500 animals in the United States, where the International Andalusian and Lusitano Horse Association (IALHA) registers around 700 new purebred foals every year. These numbers indicate that the Andalusian is a relatively rare breed in the United States. In 2003, there were 75,389 horses registered in the stud book, and they constituted almost 66 percent of the horses in Spain. Breed numbers have been increasing during the 21st century. At the end of 2010, a total of 185,926 horses were recorded in the database of the Spanish Ministerio de Medio Ambiente, y Medio Rural y Marino. Of these, 28,801 or about 15% were in other countries of the world; of those in Spain, 65,371 or about 42% were in Andalusia.
Strains and sub-types
The Carthusian Andalusian or Cartujano is generally considered the purest Andalusian strain, and has one of the oldest recorded pedigree lines in the world. The pure sub-type is rare, as only around 12 percent of the Andalusian horses registered between the founding of the stud book in the 19th century and 1998 were considered Carthusians. They made up only 3.6 percent of the overall breeding stock, but 14.2 percent of the stallions used for breeding. In the past, Carthusians were given preference in breeding, leading to a large proportion of the Andalusian population claiming ancestry from a small number of horses and possibly limiting the breed's genetic variability. A 2005 study compared the genetic distance between Carthusian and non-Carthusian horses. They calculated a Fixation index (FST) based on genealogical information and concluded that the distinction between the two is not supported by genetic evidence. However, there are slight physical differences; Carthusians have more "oriental" or concave head shapes and are more often gray in color, while non-Carthusians tend toward convex profiles and more often exhibit other coat colors such as bay.
The Carthusian line was established in the early 18th century when two Spanish brothers, Andrés and Diego Zamora, purchased a stallion named El Soldado and bred him to two mares. The mares were descended from mares purchased by the Spanish king and placed at Aranjuez, one of the oldest horse breeding farms in Spain. One of the offspring of El Soldado, a dark gray colt named Esclavo, became the foundation sire of the Carthusian line. One group of mares sired by Esclavo in about 1736 were given to a group of Carthusian monks to settle a debt. Other animals of these bloodlines were absorbed into the main Andalusian breed; the stock given to the monks was bred into a special line, known as Zamoranos. Throughout the following centuries, the Zamoranos bloodlines were guarded by the Carthusian monks, to the point of defying royal orders to introduce outside blood from the Neapolitan horse and central European breeds. They did, however, introduce Arabian and Barb blood to improve the strain. The original stock of Carthusians was greatly depleted during the Peninsular Wars, and the strain might have become extinct if not for the efforts of the Zapata family. Today, the Carthusian strain is raised in state-owned stud farms around Jerez de la Frontera, Badajoz and Cordoba, and also by several private families. Carthusian horses continue to be in demand in Spain, and buyers pay high prices for members of the strain.
Influence on other breeds
Spain's worldwide military activities between the 14th and 17th centuries called for large numbers of horses, more than could be supplied by native Spanish mares. Spanish custom also called for mounted troops to ride stallions, never mares or geldings. Due to these factors, Spanish stallions were crossed with local mares in many countries, adding Spanish bloodlines wherever they went, especially to other European breeds.
Because of the influence of the later Habsburg families, who ruled in both Spain and other nations of Europe, the Andalusian was crossbred with horses of Central Europe and the Low Countries and thus was closely related to many breeds that developed, including the Neapolitan horse, Groningen, Lipizzaner and Kladruber. Spanish horses have been used extensively in classical dressage in Germany since the 16th century. They thus influenced many German breeds, including the Hanoverian, Holstein, East Friesian and Oldenburg. Dutch breeds such as the Friesian and Gelderland also contain significant Spanish blood, as do Danish breeds such as the Frederiksborg and Knabstrupper.
Andalusians were a significant influence on the creation of the Alter Real, a strain of the Lusitano, and the Azteca, a Mexican breed created by crossing the Andalusian with American Quarter Horse and Criollo bloodlines. The Spanish jennet ancestors of the Andalusian also developed the Colonial Spanish Horse in America, which became the foundation bloodstock for many North and South American breeds. The Andalusian has also been used to create breeds more recently, with breed associations for both the Warlander (an Andalusian/Friesian cross) and the Spanish-Norman (an Andalusian/Percheron cross) being established in the 1990s.
Naming and registration
Until modern times, horse breeds throughout Europe were known primarily by the name of the region where they were bred. Thus the original term "Andalusian" simply described the horses of distinct quality that came from Andalusia in Spain. Similarly, the Lusitano, a Portuguese horse very similar to the Andalusian, takes its name from the ancient Roman province of Lusitania, now part of western Spain and most of Portugal.
The Andalusian horse has been known historically as the Iberian Saddle Horse, Iberian War Horse, Spanish Horse, Portuguese, Peninsular, Extremeño, Villanos, Zapata, Zamoranos, Castilian, and Jennet. The Portuguese name refers to what is now the Lusitano, while the Peninsular, Iberian Saddle Horse and Iberian War Horse names refer to horses from the Iberian Peninsula as a whole. The Extremeño name refers to Spanish horses from the Extremadura province of Spain and the Zapata or Zapatero name to horses that come from the Zapata family stud. The Villano name has occasionally been applied to modern Andalusians, but originally referred to heavy, crossbred horses from the mountains north of Jaen. The Carthusian horse, also known as the Carthusian-Andalusian and the Cartujano, is a sub-type of the Andalusian, rather than a distinct breed in itself. A common nickname for the Andalusian is the "Horse of Kings". Some sources state that the Andalusian and the Lusitano are genetically the same, differing only in the country of origin of individual horses.
In many areas today, the breeding, showing, and registration of the Andalusian and Lusitano are controlled by the same registries. One example of this is the International Andalusian and Lusitano Horse Association (IALHA), claimed to have the largest membership of any Andalusian registering organization. Other organizations, such as The Association of Purebred Spanish Horse Breeders of Spain ( or ANCCE), use the term or PRE to describe the true Spanish horse, and claim sole authority to officially register and issue documentation for PRE Horses, both in Spain and anywhere else in the world. In most of the world the terms "Andalusian" and "PRE" are considered one and the same breed, but the public position of the ANCCE is that terms such as "Andalusian" and "Iberian horse" refer only to crossbreds, which the ANCCE considers to be horses that lack quality and purity, without official documentation or registration from official Spanish Stud Book.
In Australasia, the Australasia Andalusian Association registers Andalusians (which the registry considers an interchangeable term for PRE), Australian Andalusians, and partbred Andalusians. They share responsibility for the Purebred Iberian Horse (an Andalusian/Lusitano cross) with the Lusitano Association of Australasia. In the Australian registry, there are various levels of crossbred horses. A first cross Andalusian is a crossbreed that is 50 percent Andalusian, while a second cross Andalusian is the result of crossing a purebred Andalusian with a first cross – resulting in a horse of 75 percent Andalusian blood. A third cross, also known by the registry as an Australian Andalusian, is when a second cross individual is mated with a foundation Andalusian mare. This sequence is known as a "breeding up" program by the registry.
Pure Spanish Horse (PRE)
The name (PRE), usually rendered in English "Pure Spanish Horse" (not a literal translation) is the term used by the ANCCE, a private organization, and the Ministry of Agriculture of Spain. The ANCCE uses neither the term "Andalusian" nor "Iberian horse", and only registers horses that have certain recognized bloodlines. In addition, all breeding stock must undergo an evaluation process. The ANCCE was founded in 1972. Spain's Ministry of Agriculture recognizes the ANCCE as the representing entity for PRE breeders and owners across the globe, as well as the administrator of the breed stud book. ANCCE functions as the international parent association for all breeders worldwide who record their horses as PRE. For example, the United States PRE association is affiliated with ANCCE, follows ANCCE rules, and has a wholly separate governance system from the IALHA.
A second group, the Foundation for the Pure Spanish Horse or PRE Mundial, has begun another PRE registry as an alternative to the ANCCE. This new registry claims that all of their registered horses trace back to the original stud book maintained by the Cria Caballar, which was a branch of the Spanish Ministry of Defense, for 100 years. Thus, the PRE Mundial registry asserts that their registry is the most authentic, purest PRE registry functioning today.
, there is a lawsuit in progress to determine the legal holder of the PRE stud book. The Unión de Criadores de Caballos Españoles (UCCE or Union of Spanish Horse Breeders) has brought a case to the highest European Union courts in Brussels, charging that the Ministry of Spain's transfer of the original PRE Libro de Origen (the official stud book) from the Cria Caballar to ANCCE was illegal. In early 2009, the courts decided on behalf of UCCE, explaining that the Cria Caballar formed the Libro de Origin. Because it was formed by a government entity, it is against European Union law for the stud book to be transferred to a private entity, a law that was broken by the transfer of the book to ANCCE, which is a non-governmental organization. The court found that by giving ANCCE sole control of the stud book, Spain's Ministry of Defense was acting in a discriminatory manner. The court held that Spain must give permission to maintain a breed stud book (called a Libro Genealógico) to any international association or Spanish national association which requests it. Based on the Brussels court decision, an application has been made by the Foundation for the Pure Spanish Horse to maintain the United States stud book for the PRE. , Spain has not revoked ANCCE's right to be the sole holder of the PRE stud book, and has instead reaffirmed the organization's status.
Uses
The Andalusian breed has over the centuries been consistently selected for athleticism. In the 17th century, referring to multi-kilometer races, Cavendish said, "They were so much faster than all other horses known at that time that none was ever seen to come close to them, even in the many remarkable races that were run." In 1831, horses at five years old were expected to be able to gallop, without changing pace, four or five leagues, about . By 1925, the Portuguese military expected horses to "cover 40 km over uneven terrain at a minimum speed of 10 km/h, and to gallop a flat course of 8 km at a minimum speed of 800 metres per minute carrying a weight of at least 70 kg", and the Spanish military had similar standards.
From the beginning of their history, Andalusians have been used for both riding and driving. Among the first horses used for classical dressage, they still compete in international competition in dressage today. At the 2002 World Equestrian Games, two Andalusians were on the bronze medal-winning Spanish dressage team, a team that went on to take the silver medal at the 2004 Summer Olympics. Today, the breed is increasingly being selectively bred for increased aptitude in classical dressage. Historically, however, they were also used as stock horses, especially suited to working with Iberian bulls, known for their aggressive temperaments. They were, and still are, known for their use in mounted bull fighting. Mares were traditionally used for la trilla, the Spanish process of threshing grain practiced until the 1960s. Mares, some pregnant or with foals at their side, spent full days trotting over the grain. As well as being a traditional farming practice, it also served as a test of endurance, hardiness and willingness for the maternal Andalusian lines.
Andalusians today are also used for show jumping, western pleasure and other horse show events. The current Traveler, the mascot of the University of Southern California, is an Andalusian. The dramatic appearance of the Andalusian horse, with its arched neck, muscular build and energetic gaits, has made it a popular breed to use in film, particularly in historical and fantasy epics. Andalusians have been present in films ranging from Gladiator to Interview with the Vampire, and Lara Croft Tomb Raider: The Cradle of Life to Braveheart. The horses have also been seen in such fantasy epics as The Lord of the Rings film trilogy, King Arthur, and The Chronicles of Narnia: The Lion, the Witch and the Wardrobe. In 2006, a rearing Andalusian stallion, ridden by Spanish conquistador Don Juan de Oñate, was recreated as the largest bronze equine in the world. Measuring high, the statue currently stands in El Paso, Texas.
| Biology and health sciences | Horses | null |
2910801 | https://en.wikipedia.org/wiki/Petroleum%20reservoir | Petroleum reservoir | A petroleum reservoir or oil and gas reservoir is a subsurface accumulation of hydrocarbons contained in porous or fractured rock formations. Such reservoirs form when kerogen (ancient plant matter) is created in surrounding rock by the presence of high heat and pressure in the Earth's crust.
Reservoirs are broadly classified as conventional and unconventional reservoirs. In conventional reservoirs, the naturally occurring hydrocarbons, such as crude oil (petroleum) or natural gas, are trapped by overlying rock formations with lower permeability, while in unconventional reservoirs the rocks have high porosity and low permeability, which keeps the hydrocarbons trapped in place, therefore not requiring a cap rock. Reservoirs are found using hydrocarbon exploration methods.
Oil field
An oil field is an area of accumulated liquid petroleum underground in multiple (potentially linked) reservoirs, trapped as it rises to impermeable rock formations. In industrial terms, an oil field implies that there is an economic benefit worthy of commercial attention. Oil fields may extend up to several hundred kilometers across the surface, meaning that extraction efforts can be large and spread out across the area. In addition to extraction equipment, there may be exploratory wells probing the edges to find more reservoir area, pipelines to transport the oil elsewhere, and support facilities.
Oil fields can occur anywhere that the geology of the underlying rock allows, meaning that certain fields can be far away from civilization, including at sea. Creating an operation at an oil field can be a logistically complex undertaking, as it involves the equipment associated with extraction and transportation, as well as infrastructure such as roads and housing for workers. This infrastructure has to be designed with the lifespan of the oil field in mind, as production can last many years. Several companies, such as Hill International, Bechtel, Esso, Weatherford International, Schlumberger, Baker Hughes and Halliburton, have organizations that specialize in the large-scale construction of the infrastructure to support oil field exploitation.
The term "oilfield" can be used as a shorthand to refer to the entire petroleum industry. However, it is more accurate to divide the oil industry into three sectors: upstream (crude oil production from wells and separation of water from oil), midstream (pipeline and tanker transport of crude oil) and downstream (refining of crude oil to products, marketing of refined products, and transportation to oil stations).
More than 65,000 oil fields are scattered around the globe, on land and offshore. The largest are the Ghawar Field in Saudi Arabia and the Burgan Field in Kuwait, with more than 66 to 104 billion barrels estimated in each. In the modern age, the location of oil fields with proven oil reserves is a key underlying factor in many geopolitical conflicts.
Gas field
Natural gas originates by the same geological thermal cracking process that converts kerogen to petroleum. As a consequence, oil and natural gas are often found together. In common usage, deposits rich in oil are known as oil fields, and deposits rich in natural gas are called natural gas fields.
In general, organic sediments buried in depths of 1,000 m to 6,000 m (at temperatures of 60 °C to 150 °C) generate oil, while sediments buried deeper and at higher temperatures generate natural gas. The deeper the source, the "drier" the gas (that is, the smaller the proportion of condensates in the gas). Because both oil and natural gas are lighter than water, they tend to rise from their sources until they either seep to the surface or are trapped by a non-permeable stratigraphic trap. They can be extracted from the trap by drilling.
The largest natural gas field is South Pars/Asalouyeh gas field, which is shared between Iran and Qatar. The second largest natural gas field is the Urengoy gas field, and the third largest is the Yamburg gas field, both in Russia.
Like oil, natural gas is often found underwater in offshore gas fields such as the North Sea, Corrib Gas Field off Ireland, and near Sable Island. The technology to extract and transport offshore natural gas is different from land-based fields. It uses a few, very large offshore drilling rigs, due to the cost and logistical difficulties in working over water.
Rising gas prices in the early 21st century encouraged drillers to revisit fields that previously were not considered economically viable. For example, in 2008 McMoran Exploration passed a drilling depth of over 32,000 feet (9754 m) (the deepest test well in the history of gas production) at the Blackbeard site in the Gulf of Mexico. ExxonMobil's drill rig there had reached 30,000 feet by 2006, without finding gas, before it abandoned the site.
Formation
Crude oil is found in all oil reservoirs formed in the Earth's crust from the remains of once-living things. Evidence indicates that millions of years of heat and pressure changed the remains of microscopic plants and animals into oil and natural gas.
Roy Nurmi, an interpretation adviser for Schlumberger oil field services company, described the process as follows:
Plankton and algae, proteins and the life that's floating in the sea, as it dies, falls to the bottom, and these organisms are going to be the source of our oil and gas. When they're buried with the accumulating sediment and reach an adequate temperature, something above 50 to 70 °C they start to cook. This transformation, this change, changes them into the liquid hydrocarbons that move and migrate, will become our oil and gas reservoir.
In addition to the aquatic ecosystem, which is usually a sea but might also be a river, lake, coral reef, or algal mat, the formation of an oil or gas reservoir also requires a sedimentary basin that passes through four steps:
Deep burial under sand and mud
Pressure cooking
Hydrocarbon migration from the source to the reservoir rock
Trapping by impermeable rock
Timing is also an important consideration; it is suggested that the Ohio River Valley could have had as much oil as the Middle East at one time, but that it escaped due to a lack of traps. The North Sea, on the other hand, endured millions of years of sea level changes that successfully resulted in the formation of more than 150 oil fields.
Although the process is generally the same, various environmental factors lead to the creation of a wide variety of reservoirs. Reservoirs exist anywhere from the land surface to below the surface and are a variety of shapes, sizes, and ages. In recent years, igneous reservoirs have become an important new field of oil exploration, especially in trachyte and basalt formations. These two types of reservoirs differ in oil content and physical properties like fracture connectivity, pore connectivity, and rock porosity.
Geology
Traps
A trap forms when the buoyancy forces driving the upward migration of hydrocarbons through a permeable rock cannot overcome the capillary forces of a sealing medium. The timing of trap formation relative to that of petroleum generation and migration is crucial to ensuring a reservoir can form.
Petroleum geologists broadly classify traps into three categories that are based on their geological characteristics: the structural trap, the stratigraphic trap, and the far less common hydrodynamic trap. The trapping mechanisms for many petroleum reservoirs have characteristics from several categories and can be known as a combination trap. Traps are described as structural traps (in deformed strata such as folds and faults) or stratigraphic traps (in areas where rock types change, such as unconformities, pinch-outs and reefs).
Structural traps
Structural traps are formed as a result of changes in the structure of the subsurface from processes such as folding and faulting, leading to the formation of domes, anticlines, and folds. Examples of this kind of trap are an anticline trap, a fault trap, and a salt dome trap. They are more easily delineated and more prospective than their stratigraphic counterparts, with the majority of the world's petroleum reserves being found in structural traps.
Stratigraphic traps
Stratigraphic traps are formed as a result of lateral and vertical variations in the thickness, texture, porosity, or lithology of the reservoir rock. Examples of this type of trap are an unconformity trap, a lens trap and a reef trap.
Hydrodynamic traps
Hydrodynamic traps are a far less common type of trap. They are caused by the differences in water pressure, that are associated with water flow, creating a tilt of the hydrocarbon-water contact.
Seal / cap rock
The seal (also referred to as a cap rock) is a fundamental part of the trap that prevents hydrocarbons from further upward migration. A capillary seal is formed when the capillary pressure across the pore throats is greater than or equal to the buoyancy pressure of the migrating hydrocarbons. They do not allow fluids to migrate across them until their integrity is disrupted, causing them to leak. There are two types of capillary seal whose classifications are based on the preferential mechanism of leaking: the hydraulic seal and the membrane seal.
A membrane seal will leak whenever the pressure differential across the seal exceeds the threshold displacement pressure, allowing fluids to migrate through the pore spaces in the seal. It will leak just enough to bring the pressure differential below that of the displacement pressure and will reseal.
A hydraulic seal occurs in rocks that have a significantly higher displacement pressure such that the pressure required for tension fracturing is actually lower than the pressure required for fluid displacement—for example, in evaporites or very tight shales. The rock will fracture when the pore pressure is greater than both its minimum stress and its tensile strength then reseal when the pressure reduces and the fractures close.
Unconventional reservoirs
Unconventional (oil & gas) reservoirs are accumulations where oil and gas phases are tightly bound to the rock fabric by strong capillary forces, requiring specialised measures for evaluation and extraction. Unconventional reservoirs form in completely different ways to conventional reservoirs, the main difference being that they do not have "traps". This type of reservoir can be driven in a unique way as well, as buoyancy might not be the driving force for oil and gas accumulation in such reservoirs. This is analogous to saying that the oil which can be extracted forms within the source rock itself, as opposed to accumulating under a cap rock. Oil sands are an example of an unconventional oil reservoir.
Unconventional reservoirs and their associated unconventional oil encompass a broad spectrum of petroleum extraction and refinement techniques, as well as many different sources. Since the oil is contained within the source rock, unconventional reservoirs require that the extracting entity function as a mining operation rather than drilling and pumping like a conventional reservoir. This has tradeoffs, with higher post-production costs associated with complete and clean extraction of oil being a factor of consideration for a company interested in pursuing a reservoir. Tailings are also left behind, increasing cleanup costs. Despite these tradeoffs, unconventional oil is being pursued at a higher rate because of the scarcity of conventional reservoirs around the world.
Estimating reserves
After the discovery of a reservoir, a petroleum engineer will seek to build a better picture of the accumulation. In a simple textbook example of a uniform reservoir, the first stage is to conduct a seismic survey to determine the possible size of the trap. Appraisal wells can be used to determine the location of oil-water contact and with it the height of the oil bearing sands. Often coupled with seismic data, it is possible to estimate the volume of an oil-bearing reservoir.
The next step is to use information from appraisal wells to estimate the porosity of the rock. The porosity of an oil field, or the percentage of the total volume that contains fluids rather than solid rock, is 20–35% or less. It can give information on the actual capacity. Laboratory testing can determine the characteristics of the reservoir fluids, particularly the expansion factor of the oil, or how much the oil expands when brought from the high pressure and high temperature of the reservoir to a "stock tank" at the surface.
With such information, it is possible to estimate how many "stock tank" barrels of oil are located in the reservoir. Such oil is called the stock tank oil initially in place. As a result of studying factors such as the permeability of the rock (how easily fluids can flow through the rock) and possible drive mechanisms, it is possible to estimate the recovery factor, or what proportion of oil in place can be reasonably expected to be produced. The recovery factor is commonly 30–35%, giving a value for the recoverable resources.
The difficulty is that reservoirs are not uniform. They have variable porosities and permeabilities and may be compartmentalized, with fractures and faults breaking them up and complicating fluid flow. For this reason, computer modeling of economically viable reservoirs is often carried out. Geologists, geophysicists, and reservoir engineers work together to build a model that allows simulation of the flow of fluids in the reservoir, leading to an improved estimate of the recoverable resources.
Reserves are only the part of those recoverable resources that will be developed through identified and approved development projects. Because the evaluation of reserves has a direct impact on the company or the asset value, it usually follows a strict set of rules or guidelines.
Production
To obtain the contents of the oil reservoir, it is usually necessary to drill into the Earth's crust, although surface oil seeps exist in some parts of the world, such as the La Brea Tar Pits in California and numerous seeps in Trinidad. Factors that affect the quantity of recoverable hydrocarbons in a reservoir include the fluid distribution in the reservoir, initial volumes of fluids in place, reservoir pressure, fluid and rock properties, reservoir geometry, well type, well count, well placement, development concept, and operating philosophy.
Modern production includes thermal, gas injection, and chemical methods of extraction to enhance oil recovery.
Drive mechanisms
A virgin reservoir may be under sufficient pressure to push hydrocarbons to the surface. As the fluids are produced, the pressure will often decline, and production will falter. The reservoir may respond to the withdrawal of fluid in a way that tends to maintain the pressure. Artificial drive methods may be necessary.
Solution-gas drive
This mechanism (also known as depletion drive) depends on the associated gas of the oil. The virgin reservoir may be entirely semi-liquid but will be expected to have gaseous hydrocarbons in solution due to the pressure. As the reservoir depletes, the pressure falls below the bubble point, and the gas comes out of solution to form a gas cap at the top. This gas cap pushes down on the liquid helping to maintain pressure.
This occurs when the natural gas is in a cap below the oil. When the well is drilled the lowered pressure above means that the oil expands. As the pressure is reduced it reaches bubble point, and subsequently the gas bubbles drive the oil to the surface. The bubbles then reach critical saturation and flow together as a single gas phase. Beyond this point and below this pressure, the gas phase flows out more rapidly than the oil because of its lowered viscosity. More free gas is produced, and eventually the energy source is depleted. In some cases depending on the geology the gas may migrate to the top of the oil and form a secondary gas cap. Some energy may be supplied by water, gas in water, or compressed rock. These are usually minor contributions with respect to hydrocarbon expansion.
By properly managing the production rates, greater benefits can be had from solution-gas drives. Secondary recovery involves the injection of gas or water to maintain reservoir pressure. The gas/oil ratio and the oil production rate are stable until the reservoir pressure drops below the bubble point when critical gas saturation is reached. When the gas is exhausted, the gas/oil ratio and the oil rate drops, the reservoir pressure has been reduced, and the reservoir energy is exhausted.
Gas cap drive
In reservoirs already having a gas cap (the virgin pressure is already below bubble point), the gas cap expands with the depletion of the reservoir, pushing down on the liquid sections applying extra pressure. This is present in the reservoir if there is more gas than can be dissolved in the reservoir. The gas will often migrate to the crest of the structure. It is compressed on top of the oil reserve, as the oil is produced the cap helps to push the oil out. Over time the gas cap moves down and infiltrates the oil, and the well will produce more and more gas until it produces only gas.
It is best to manage the gas cap effectively, that is, placing the oil wells such that the gas cap will not reach them until the maximum amount of oil is produced. Also a high production rate may cause the gas to migrate downward into the production interval. In this case, over time the reservoir pressure depletion is not as steep as in the case of solution-based gas drive. In this case, the oil rate will not decline as steeply but will depend also on the placement of the well with respect to the gas cap. As with other drive mechanisms, water or gas injection can be used to maintain reservoir pressure. When a gas cap is coupled with water influx, the recovery mechanism can be highly efficient.
Aquifer (water) drive
Water (usually salty) may be present below the hydrocarbons. Water, as with all liquids, is compressible to a small degree. As the hydrocarbons are depleted, the reduction in pressure in the reservoir allows the water to expand slightly. Although this unit expansion is minute, if the aquifer is large enough this will translate into a large increase in volume, which will push up on the hydrocarbons, maintaining pressure.
With a water-drive reservoir, the decline in reservoir pressure is very slight; in some cases, the reservoir pressure may remain unchanged. The gas/oil ratio also remains stable. The oil rate will remain fairly stable until the water reaches the well. In time, the water cut will increase, and the well will be watered out.
The water may be present in an aquifer (but rarely one replenished with surface water). This water gradually replaces the volume of oil and gas that is produced out of the well, given that the production rate is equivalent to the aquifer activity. That is, the aquifer is being replenished from some natural water influx. If the water begins to be produced along with the oil, the recovery rate may become uneconomical owing to the higher lifting and water disposal costs.
Water and gas injection
If the natural drives are insufficient, as they very often are, then the pressure can be artificially maintained by injecting water into the aquifer or gas into the gas cap.
Gravity drainage
The force of gravity will cause the oil to move downward of the gas and upward of the water. If vertical permeability exists then recovery rates may be even better.
Gas and gas condensate reservoirs
These occur if the reservoir conditions allow the hydrocarbons to exist as a gas. Retrieval is a matter of gas expansion. Recovery from a closed reservoir (i.e., no water drive) is very good, especially if bottom hole pressure is reduced to a minimum (usually done with compressors at the wellhead). Any produced liquids are light-colored to colorless, with a gravity higher than 45 API. Gas cycling is the process where dry gas is injected and produced along with condensed liquid.
| Technology | Fuel | null |
2911772 | https://en.wikipedia.org/wiki/Chalicotheriidae | Chalicotheriidae | Chalicotheriidae (from Greek chalix, "gravel" and therion, "beast") is an extinct family of herbivorous, odd-toed ungulate (perissodactyl) mammals that lived in North America, Eurasia, and Africa from the Middle Eocene to the Early Pleistocene. They are often called chalicotheres, a term which is also applied to the broader grouping of Chalicotherioidea. They are noted for their unusual morphology compared to other ungulates, such as their clawed forelimbs. Members of the subfamily Chalicotheriinae developed elongate gorilla-like forelimbs that are thought to have been used to grasp vegetation. They are thought to have been browsers on foliage as well as possibly bark and fruit.
History of discovery
The earliest remains chalicotheres were discovered were ungual phalanges found near Eppelsheim, Germany in early 19th century. These remains were considered to belong to gigantic pangolins by Georges Cuvier in 1822 while Johann Jakob Kaup in 1833 alternatively attributed them to deinotheres (which at that time were considered enigmatic). Also in 1833, Kaup described chalicothere teeth as belonging to the new genus Chalicotherium from the same locality, which he did not recognise as belonging to the same species as the ungual phalanges. Beginning in the 1930s, the Sansan deposit in southern France was excavated for fossils, yielding remains of chalicotheres. In 1837, postcranial remains from the deposit were given the name "Macrotherium" by Édouard Lartet, who like Cuvier thought the remains represented those of a giant pangolin/edentate. In 1849, Henri Marie Ducrotay de Blainville described chalicothere skull remains from Sansan as belonging to the ungulate genus Anoplotherium (an extinct even-toed ungulate now known to be unrelated to chalicotheres). It was only in 1890 that a complete chalicothere skeleton found at Sansan was described by Henri Filhol, showing the skulls/teeth and the postcranial remains belonged to the same unusual animal.
Description
Unlike modern perissodactyls, chalicotheres had clawed feet. They had, lower incisors that cropped food against a toothless pad in the upper jaw, low-crowned molar teeth, and were browsers on trees and shrubs throughout their history. They evolved in two different directions, which became separate subfamilies, the Schizotheriinae and the Chalicotheriinae.
Schizotherine chalicotheres such as Moropus had relatively equal length limbs, and lived in a variety of forest, woodland, and savannah habitats in Asia, Africa, and North America. They developed long necks and skull adaptations that suggest they had long, extensible tongues to reach browse, like those of giraffes. Strong hindlimbs and an elongated pelvis suggest they could have reared upright as modern goats do, and used their front claws to pull branches within reach of the tongue. The claws were retractable, and they walked normally on the bottom of the foot. Studies of tooth wear suggest they ate leaves, twigs, fruit, and bark.
Chalicotheriines, such as Anisodon, lived only in moist, closed-canopy forests, never reached the Americas, and developed very unusual anatomy for an ungulate. Their shorter necks and horse-like heads did not show adaptations to reach high. Instead, they developed very long forelimbs with mobile shoulder joints and hooklike claws. The pelvis and hindlimbs were specialized to stand upright, and to sit for hours while feeding, like the living gelada monkey. Some early paleontologists thought the claws were used to dig up roots and tubers, but their teeth were designed for soft foods, and studies of tooth wear show they ate fruit and seeds. Their forelimbs were specialized to reach, grasp, and strip or sweep plants to the mouth. They could not retract the huge front claws, and knuckle-walked on their forelimbs. The chalicotheriines' anatomical design, posture, and locomotion show convergence with other large browsers that feed selectively in a bipedal position, such as the ground sloths, gorillas, and giant pandas.
Chalicothere fossils are uncommon even in areas where other taxa of similar size are well-preserved, which suggests they were mostly solitary animals, and unlike horses, rhinos, and brontotheres, never evolved species that lived in herds. Only two species of chalicothere are known from complete skeletons, the schizotheriine Moropus from the early Miocene of North America, and the chalicotheriine Anisodon from the middle Miocene of Europe. Fossils of other species range from very fragmentary to moderately complete. Chalicotheres ranged in size from an antelope to a large draft horse.
Evolution
Chalicotheres are part of the order Perissodactyla, which includes modern equines, rhinoceroses, and tapirs, as well as extinct relatives like brontotheres. As the early evolution of perissodactyls is still unresolved, their closest relatives among other perissodactyl groups is obscure. They are generally placed as part of the clade Ancylopoda alongside their close relatives Lophiodontidae. Many studies considered them as closer to Ceratomorpha (which includes tapirs and rhinoceroses) than Equoidea. A 2004 cladistic study alternatively recovered Ancylopoda as sister to all modern perissodactyls (which includes Equoidea and Ceratomorpha), with the brontotheres as the most distantly related within the order Perissodactyla.
Chalicotheres can be first identified with certainty around 46 million years ago, in the Eocene of Asia. The family is thought to have evolved there, but appeared in North America by the Eocene. By the late Oligocene, they had divided into schizotheriines and chalicotheriines. (Earlier chalicotheres are often referred to the family Eomoropidae; it is not yet clear whether they had claws or how the two subfamilies diverged.) Both subfamilies were successful over many millions of years, and reached their greatest diversity in the Miocene. Advanced schizotheriines (Moropus) entered North America via the Bering land bridge at the Oligicene-Miocene boundary, and expanded southward into Central America.
Never common animals, the chalicotheres declined from the late Neogene onwards, disappearing from North America and Europe by end of the Miocene. The youngest chalicotheres are the chalicotheriines Hesperotherium from the Early Pleistocene of China, Nestoritherium from the Early Pleistocene of Myanmar, as well as the schizotheriine Ancylotherium from the Early Pleistocene of Eastern and Southern Africa, also possibly known from the Early Pleistocene of China.
| Biology and health sciences | Perissodactyla | Animals |
2915941 | https://en.wikipedia.org/wiki/Pouteria%20sapota | Pouteria sapota | Pouteria sapota, the mamey sapote, is a species of tree native to Mexico and Central America. The tree is also cultivated in the Caribbean. Its fruit is eaten in many Latin American countries. The fruit is made into foods such as milkshakes and ice cream.
Some of its names in Latin American countries, such as (Cuba), (Costa Rica) and (South America), refer to the reddish colour of its flesh to distinguish it from the unrelated but similar-looking Mammea americana, whose fruit is usually called "yellow mamey" ().
The Australian and Queensland governments' research and development programs have grown mamey sapote in Australia.
Description
Mamey sapote is a large and highly ornamental evergreen tree that can reach a height of at maturity. It is mainly propagated by grafting, which ensures the new plant has the same characteristics as the parent, especially its fruit, as it does not grow true to seed. It is also considerably faster than growing trees by seed, producing fruit in 3–5 years; trees grown from seed require 7 years of growth before fruiting. In Florida, the fruit is harvested from May to July with some cultivars available all year.
The fruit, technically a berry, is about long and wide and has flesh ranging in color from pink to orange to red. The brown skin has a texture somewhat between sandpaper and the fuzz on a peach. The fruit's texture is creamy and soft, and the flavor is a mix of sweet potato, pumpkin, honey, prune, peach, apricot, cantaloupe, cherry, and almond. A mamey sapote is ripe when the flesh is vibrant salmon in color when a fleck of the skin is removed. The flesh should give slightly, as with an overripe avocado. The leaves are pointed at both ends, 4 to 12 inches in length, and grow in clusters at the ends of branches.
The mamey sapote is related to other sapotes such as sapodilla (Manilkara zapota), abiu (P. caimito), and canistel (P. campechiana), but unrelated to the black sapote (Diospyros digyna) and white sapote (Casimiroa edulis).
Uses
The fruit is eaten raw or made into milkshakes, smoothies, ice cream, and fruit bars. It can be used to produce marmalade and jelly. It can also be fried like bacon. Some beauty products use oil pressed from the seed, otherwise known as sapayul oil.
Nutrition
The fruit is an excellent source of vitamins B6 and C, and is a good source of riboflavin, niacin, vitamin E, manganese, potassium, and dietary fiber. Research has identified several new carotenoids from the ripe fruit.
Synonyms
| Biology and health sciences | Tropical and tropical-like fruit | Plants |
33183231 | https://en.wikipedia.org/wiki/Freshwater%20crab | Freshwater crab | Around 1,300 species of freshwater crabs are distributed throughout the tropics and subtropics, divided among eight families. They show direct development and maternal care of a small number of offspring, in contrast to marine crabs, which release thousands of planktonic larvae. This limits the dispersal abilities of freshwater crabs, so they tend to be endemic to small areas. As a result, a large proportion are threatened with extinction.
Systematics
More than 1,300 described species of freshwater crabs are known, out of a total of 6,700 species of crabs across all environments. The total number of species of freshwater crabs, including undescribed species, is thought to be up to 65% higher, potentially up to 2,155 species, although most of the additional species are currently unknown to science. They belong to eight families, each with a limited distribution, although various crabs from other families are also able to tolerate freshwater conditions (euryhaline) or are secondarily adapted to fresh water. The phylogenetic relationships between these families is still a matter of debate, so how many times the freshwater lifestyle has evolved among the true crabs is unknown. The eight families are:
Superfamily Trichodactyloidea
Trichodactylidae (Central America and South America)
Superfamily Potamoidea
Potamidae (Mediterranean Basin and Asia)
Potamonautidae (Africa, including Madagascar)
Deckeniidae (East Africa and Seychelles) – also treated as part of Potamonautidae
Platythelphusidae (East Africa) – also treated as part of Potamonautidae
Superfamily Gecarcinucoidea
Gecarcinucidae (Asia)
Parathelphusidae (Asia and Australasia) – nowadays treated as a junior synonym of Gecarcinucidae
Superfamily Pseudothelphusoidea
Pseudothelphusidae (Central America and South America)
The fossil record of freshwater organisms is typically poor, so few fossils of freshwater crabs have been found. The oldest is Tanzanonautes tuerkayi, from the Oligocene of East Africa, and the evolution of freshwater crabs is likely to postdate the break-up of the supercontinent Gondwana.
Members of the family Aeglidae and Clibanarius fonticola are also restricted to fresh water, but these "crab-like" crustaceans are members of the infraorder Anomura (true crabs are Brachyura).
Description and lifecycle
The external morphology of freshwater crabs varies very little, so the form of the gonopod (first abdominal appendage, modified for insemination) is of critical importance for classification. Development of freshwater crabs is characteristically direct, where the eggs hatch as juveniles, with the larval stages passing within the egg. The broods comprise only a few hundred eggs (compared to hundreds of thousands for marine crabs), each of which is quite large, at a diameter around .
The colonisation of fresh water has required crabs to alter their water balance; freshwater crabs can reabsorb salt from their urine, and have various adaptations to reduce the loss of water. In addition to their gills, freshwater crabs have a "pseudolung" in their gill chamber that allows them to breathe in air. These developments have preadapted freshwater crabs for terrestrial living, although freshwater crabs need to return to water periodically to excrete ammonia.
Ecology and conservation
Freshwater crabs are found throughout the tropical and subtropical regions of the world. They live in a wide range of water bodies, from fast-flowing rivers to swamps, as well as in tree boles or caves. They are primarily nocturnal, emerging to feed at night; most are omnivores, although a small number are specialist predators, such as Platythelphusa armata from Lake Tanganyika, which feeds almost entirely on snails. Some species provide important food sources for various vertebrates. A number of freshwater crabs (for example species from the genus Nanhaipotamon) are secondary hosts of flukes in the genus Paragonimus, which causes paragonimiasis in humans.
The majority of species are narrow endemics, occurring in only a small geographical area. This is at least partly attributable to their poor dispersal abilities and low fecundity, and to habitat fragmentation caused by the world's human population. In West Africa, species that live in savannas have wider ranges than species from the rainforest; in East Africa, species from the mountains have restricted distributions, while lowland species are more widespread.
Every species of freshwater crab described so far has been assessed by the International Union for Conservation of Nature; of the species for which data are available, 32% are threatened with extinction. For instance, all but one of Sri Lanka's 50 freshwater crab species are endemic to that country, and more than half are critically endangered.
| Biology and health sciences | Crabs and hermit crabs | Animals |
5264689 | https://en.wikipedia.org/wiki/Yukon%20Gold%20potato | Yukon Gold potato | 'Yukon Gold' is a large cultivar of potato most distinctly characterized by its thin, smooth, eye-free skin and yellow-tinged flesh. This potato was developed in the 1960s by Garnet ("Gary") Johnston in Guelph, Ontario, Canada, with the help of Geoff Rowberry at the University of Guelph. The official cross bred strain was made in 1966 and 'Yukon Gold' was finally released into the market in 1980.
Development and naming
In the 1900s, many Dutch and Belgian immigrants began settling in the "Banana Belt" region of southern Ontario. Many of these immigrants began vegetable farming around the towns of Simcoe, Leamington, and Harrow along the shore of Lake Erie. In the 1950s, the vegetable growers of this region began petitioning for the breeding rights and licensing for a yellow-fleshed potato like ones they were used to growing in Europe. For Gary Johnston, this began the nearly 30-year development of the 'Yukon Gold' potato.
In 1953, Johnston was a lab technician in the potato development laboratory at the Ontario Agriculture College and he led a team that cross-bred two varieties to create the new type. In 1959, one of Johnston’s graduate students, a young man originally from Peru, told him of a small, rough, deep-yellow-fleshed potato (Solanum goniocalyx, known as , Spanish for "yellow potato") that was grown by the many indigenous communities in the Peruvian Andes. In Lima, this cultivar is considered a delicacy for its bright colour and distinct flavour. After trying these Peruvian potatoes, Johnston set out to breed a potato with the same colour and flavor characteristics, but larger in size and with a smoother shape, similar to the potatoes being grown in that part of Southwestern Ontario. In 1966, the development team made their first cross between a W5289-4 (2× cross between 'Yema de huevo' and 2× Katahdin) and a 'Norgleam' potato native to North Dakota. After the 66th cross that year, a true-breeding seed was produced, and the G6666 was created.
The early name for the new cultivar was "Yukon", for the Yukon River involved in the Klondike Gold Rush in Northern Canada. Charlie Bishop, or Walter Shy according to some sources, suggested adding "Gold" to describe the colour and appearance. "It was a revolutionary concept ... He was a pioneer. He [Johnston] had the vision for yellow-fleshed potatoes", said Hielke De Jong, a potato breeder with Agriculture and Agri-Food Canada. Johnston also developed and brought 15 other potato varieties to market while at the Ontario Agriculture College lab, where he had been seconded by his employer, Agriculture Canada.
A University publication states that "Yukon Gold was the first Canadian-bred potato variety to be promoted, packaged and marketed with its name right on the pack".
In spite of the overwhelming success of this potato for some years, sales in Canada dropped 30% between 2004 and 2014 as other varieties became increasingly popular.
Yukon Gold potatoes are susceptible to seed decay, blackleg, early blight, late blight, early dying, PVY, soft rot, dry rot, leak, pink rot, silver scurf, and black scurf.
Storage
This cultivar is resistant to bruising and does not sprout a lot because it has good dormancy. If tubers are stored correctly they will not lose a lot of moisture compared to other cultivars. It is important that the lenticels are not swollen and that the skin is not bruised because this can lead to major rot issues.
| Biology and health sciences | Root vegetables | Plants |
5269284 | https://en.wikipedia.org/wiki/Altostratus%20undulatus%20cloud | Altostratus undulatus cloud | The altostratus undulatus is a type of altostratus cloud with signature undulations within it. These undulations may be visible (usually as "wavy bases"), but frequently they are indiscernible to the naked eye. These formations will generally appear in the early stages of destabilizing return flows, especially over the southern plains of the United States, when the surface temperature is still relatively cool. The wavy strips of clouds are generally near an inversion surface.
Also referred to as billow clouds, wind row clouds, or wave clouds, variations of the undulatus can be elements that have merged or single elements that have stretched through the sky.
| Physical sciences | Clouds | Earth science |
38662000 | https://en.wikipedia.org/wiki/Port%20of%20Port%20Hedland | Port of Port Hedland | Port Hedland is one of the largest iron ore loading ports in the world and the largest in Australia. In 2022, it had the largest bulk cargo throughput in Australia. With the neighbouring ports of Port Walcott and Dampier, Port Hedland is one of three major iron ore exporting ports in the Pilbara region of Western Australia.
History
Named after Captain Hedland, the Master of a ship that anchored there in 1863, Port Hedland was first developed in order to service the needs of the local pastoral industry in East Pilbara. The first jetty was built in 1896, this was extended in 1908 after the discovery of gold in the Marble Bar area.
Until the 1930s the port was predominantly used to import goods and stores for the local industries and to export pearl, shell, wool, livestock, gold, tin and copper.
With the end of World War II, the port began exporting significant amounts of manganese.
The 1960s saw the development of the port by the iron ore and salt industries. Mount Goldsworthy Mining Associates, a company later absorbed by BHP, dredged an approach channel and turning basin for 65,000 DWT ships. Meanwhile, the Leslie Salt Company, from August 2001 Dampier Salt (part of Rio Tinto), built a land backed wharf and facilities to aid salt exports and fuel imports.
Further dredging was performed after the Mount Newman Mining Company, a subsidiary of BHP, chose Port Hedland as its export port. The new works allowed for ships up to 120,000 DWT.
Between the 1960s and today and extensive dredging and building has taken Port Hedland from a convenient anchorage to 15 berths capable of loading various ores and goods onto ships ranging from 25,000 DWT to 320,000 DWT.
In 2005/06 Port Hedland became the first Australian port to export in excess of 100 million tonnes annually. In 2010/11 the port exported a record 199 million tonnes, making it the largest port by cargo tonnage in Australia.
In 2021/22, 561 million tonnes of cargo passed through the port.
Port authority
Port Hedland's harbour is managed by the Pilbara Port Authority, a Government of Western Australia instrumentality. The port authority's headquarters, control tower and heliport are at Mangrove Point, just to the west of The Esplanade at the western end of Port Hedland.
The tugboat pen, customs office and public jetty are at nearby Laurentius Point. The harbour's wharves are located on both sides of the harbour – Finucane Island to the west and Port Hedland to the east.
Geography
The points are key references to the port and its composition. Mangrove and Laurentius points are already mentioned, Utah Point, and Anderson Point are located in the inner part of the harbour and exist within the complex of current berths.
Berths
Allocation of berths includes a commercial in confidence agreement between the PPA and BHP known as the Harriet Point agreement
Finucane Island berths - BHP berths FIA, FIB, FIC & FID, PPA Wharf 4
Anderson Point berths - Fortescue Metals Group berths 1,2,3 and 4,5 (in South West Creek)
South West Creek - Roy Hill Stanley Point SP 1 & SP 2
Inner Harbour berths - PPA Wharves 1,2 & 3 and BHP berths Nelson Point NPA, NPB, NPC, NPD
Gallery
Access by oceangoing vessels into and out of the harbour is via a narrow curved channel. The following series of images depicts a 225 m (246.1 yd) long bulk carrier, Darya Shanthi, using the channel to enter the harbour. Visible in the foreground of each image is part of the harbour's system of mangroves.
Port statistics
| Technology | Specific piers and ports | null |
1446187 | https://en.wikipedia.org/wiki/Lime%20%28color%29 | Lime (color) | Lime is a color that is a shade of yellow-green, so named because it is a representation of the color of the citrus fruit called limes. It is the color that is in between the web color chartreuse and yellow on the color wheel. Alternate names for this color included yellow-green, lemon-lime, lime green, or bitter lime.
The first recorded use of lime green as a color name in English was in 1890.
Lime (color hex code #C0FF00) is a pure spectral color at approximately 564 nanometers on the visible spectrum when plotted on the CIE chromaticity diagram.
Usage
During the 2000s, lime green was a very popular aesthetic, particularly with products, throughout the entirety of the decade and eventually saw a resurgence during the early 2020s. Famous examples include Song Airlines, Crocs shoes, and the Seattle Seahawks.
Some fire engines in the United States are lime yellow rather than red due to safety and ergonomics reasons. A 2009 study by the U.S. Fire Administration concluded that fluorescent colors, including yellow-green and orange, are easiest to spot in daylight.
The National Rugby League team Canberra Raiders uses lime green as one of its main colors, as does the National Football League's Seattle Seahawks, which utilizes a color officially named Action Green, which strongly resembles lime green.
Variations
Key lime
Key lime is a light lime color that is named after a Crayola Pearl Brites crayon.
Lemon-lime
Lemon-lime is a fluorescent chartreuse color that is named after carbonated soft drinks such as Sprite, 7 Up, and Sierra Mist. The red value to this neon color is almost to yellow.
Arctic lime
The color Arctic lime is close to electric lime, and was named in 2009. This is one of the colors in Crayola's eXtreme colors ultra-bright colored pencils.
Peridot
The color peridot is a shade of lime with lemon undertones, which represents the color of the peridot gemstone. Peridot is the birthstone for those born in August.
Volt
The color Volt is used by Nike in several of their athletic products, most notably their Air Max 90 Hyperfuse sneakers, which were introduced in 2011. This color is similar to electric lime.
Electric lime
Electric lime is a Crayola color created in 1990. This tint of lime is popular in psychedelic art.
French lime
The color French lime is the shade of lime called "lime" in the Pourpre.com color list, a color list widely popular in France.
Web color "lime" (X11 Green)
The web color named "lime", in the CSS color scheme maintained by the World Wide Web Consortium (W3C), has the identical normalized color coordinates as the color green, as found in X11 color names formulized over 1985–1989. The web color lime / X11 color green match the green primary color of the RGB color model.
The W3C web color named green is darker than the color named green in X11, using the HTML color code #008000 as compared to the color code #00FF00 in X11. This lime versus green issue is one of the very few clashes between web and X11 colors in the CSS color scheme.
Lime green
Lime green is a vivid yellowish green web color.
Bright lime
Bright lime is a luminous vivid chartreuse green web color.
| Physical sciences | Colors | Physics |
1447448 | https://en.wikipedia.org/wiki/Key%20lime | Key lime | The Key lime or acid lime (Citrus × aurantiifolia or C. aurantifolia) is a citrus hybrid (C. hystrix × C. medica) native to tropical Southeast Asia. It has a spherical fruit, in diameter. The Key lime is usually picked while it is still green, but it becomes yellow when ripe.
The Key lime has thinner rind and is smaller, seedier, more acidic and more aromatic than the Persian lime (Citrus × latifolia). It is valued for its characteristic flavor. The name comes from its association with the Florida Keys, where it is best known as the flavoring ingredient in Key lime pie. Not to be confused with West Indian lime, bartender's lime, Omani lime, or Mexican lime which are slightly different. The last classified as a distinct race with a thicker skin and darker green colour. Philippine varieties have various names, including dayap and bilolo.
Etymology
The English word lime was derived, via Spanish then French, from the Arabic word līma, which is, in turn, a derivation of the Persian word limu . Key is from Florida Keys, where the fruit was naturalised. The earliest known use of the name is from 1905, where the fruit was described as "the finest on the market. It is aromatic, juicy, and highly superior to the lemon."
Description
C. aurantiifolia is a shrubby tree, growing to , with many thorns. Dwarf varieties exist that can be grown indoors during winter months and in colder climates. Its trunk, which rarely grows straight, has many branches, and they often originate quite far down on the trunk. The leaves are ovate, long, resembling orange leaves (the scientific name aurantiifolia referring to this resemblance to the leaves of Citrus aurantium). The flowers are in diameter, are yellowish white with a light purple tinge on the margins. Flowers and fruit appear throughout the year, but are most abundant from May to September in the Northern Hemisphere.
Skin contact can sometimes cause phytophotodermatitis, which makes the skin especially sensitive to ultraviolet light.
Taxonomy
The Key lime cultivar is a citrus hybrid, Citrus micrantha × Citrus medica (a papeda-citron cross).
The Key lime has given rise to several other lime varieties. The best known, the triploid progeny of a Key lime-lemon cross, is the Persian lime (Citrus × latifolia), the most widely produced lime, globally. Others are, like their parent, classed within C. aurantiifolia. Backcrossing with citron has produced a distinct group of triploid limes that are also of commercial value to a limited degree, the seedy Tanepeo, Coppenrath, Ambilobe and Mohtasseb lime varieties as well as the Madagascar lemon. Hybridization with a mandarin-pomelo cross similar to the oranges has produced the Kirk lime. The New Caledonia and Kaghzi limes appear to have resulted from an F2 Key lime self-pollination, while a spontaneous genomic duplication gave us the tetraploid Giant Key lime. The potential to produce a wider variety of lime hybrids from the Key lime due to its tendency to form diploid gametes may reduce the disease risk presented by the limited diversity of the current commercial limes.
Distribution and habitat
C. aurantiifolia is native to Southeast Asia. Its apparent path of introduction was through the Middle East to North Africa, then to Sicily and Andalucia and then, via Spanish explorers, to the West Indies, including the Florida Keys. Henry Perrine is credited with introducing the Key lime to Florida. From the Caribbean, lime cultivation spread to tropical and subtropical North America, including Mexico, Florida, and later California.
Cultivation
History
In California in the late 19th century, "Mexican" limes were more highly valued than lemons; however, in Florida, they were generally considered weeds. Then, in 1894–95, the Great Freeze destroyed the Florida lemon groves, and farmers replanted Mexican limes instead; they soon became known as the Florida Key Lime, a "beloved regional crop". But when the 1926 Miami hurricane ripped them up, they were replanted with the hardier, thornless Persian limes.
Since the North American Free Trade Agreement came into effect, most Key limes on the United States market have been grown in Mexico, Central America and South America. They are also grown in Texas, Florida, and California.
Propagation
There are various approaches to the cultivation of Key limes. This variety of citrus can be propagated from seed and will grow true to the parent. The seeds must be kept moist until they can be planted, as they will not germinate if allowed to dry out. If the plants are propagated from seed, the seeds should be stored at least 5–6 months before planting. Alternatively, vegetative propagation from cuttings or by air layering may permit fruit production within one year, and from genetically more predictable lines of plants. Another method, digging around a mature tree to sever roots, will encourage new sprouts that can be transplanted to another location. Clones are often bud grafted into rough lemon or bitter orange to obtain strong root stocks.
It is often advisable to graft the plants onto rootstocks with low susceptibility to gummosis because seedlings generally are highly vulnerable to the disease. Useful rootstocks include wild grapefruit, cleopatra mandarin and tahiti limes. C. macrophylla is also sometimes used as a rootstock in Florida to add vigor.
Climatic conditions and fruit maturation are crucial in cultivation of the lime tree. Under consistently warm conditions potted trees can be planted at any season, whereas in cooler temperate regions it is best to wait for the late winter or early spring. The Key lime tree does best in sunny sites, well-drained soils, good air circulation, and protection from cold wind. Because its root system is shallow, the Key lime is planted in trenches or into prepared and broken rocky soil to give the roots a better anchorage and improve the trees' wind resistance. Pruning and topping should be planned to maximise the circulation of air and provide plenty of sunlight. This keeps the crown healthily dry, improves accessibility for harvesting, and discourages the organisms that cause gummosis.
Harvesting
The method of cultivation greatly affects the size and quality of the harvest. Trees cultivated from seedlings take 4–8 years before producing a harvest. They attain their maximal yield at about 10 years of age. Trees produced from cuttings and air layering bear fruit much sooner, sometimes producing fruit (though not a serious harvest) a year after planting. It takes approximately 9 months from the blossom to the fruit. When the fruit have grown to harvesting size and begin to turn yellow they are picked and not clipped. To achieve produce of the highest market value, it is important not to pick the fruit too early in the morning; the turgor is high then, and handling turgid fruit releases the peel oils and may cause spoilage.
Postharvest process
Shelf life of Key limes is an important consideration in marketing. The lime still ripens for a considerable time after harvesting, and it is usually stored between at a relative humidity of 75–85%. Special procedures are employed to control the shelf life; for example, applications of growth regulators, fruit wax, fungicides, precise cooling, calcium compounds, silver nitrate, and special packing material. The preferred storage conditions are temperatures of and a humidity over 85%, but even in ideal conditions post-harvesting losses are high.
In India most Key lime producers are small-scale farmers without access to such post-harvesting facilities, but makeshift expedients can be of value. One successful procedure is a coating of coconut oil that improves shelf life, thereby achieving a constant market supply of Key limes.
Key limes are made into black lime by boiling them in brine and drying them. Black lime is a condiment commonly used in the Middle East.
Yield
The yield varies depending on the age of the trees. Five- to seven-year-old orchards may yield about 6 t/ha (2.7 tons/acre), with harvests increasing progressively until they stabilise at about 12–18 t/ha (5.4–8 tons/acre). Seedling trees take longer to attain their maximal harvest, but eventually out-yield grafted trees.
In culture
The annual Key Lime Festival in Key West, Florida, has been held every year since 2002 over the Independence Day weekend and is a celebration of Key limes in food, drinks, and culture, with a significant emphasis on Key lime pie.
| Biology and health sciences | Citrus fruits | Plants |
1447904 | https://en.wikipedia.org/wiki/Kerr%E2%80%93Newman%20metric | Kerr–Newman metric | The Kerr–Newman metric describes the spacetime geometry around a mass which is electrically charged and rotating. It is a vacuum solution which generalizes the Kerr metric (which describes an uncharged, rotating mass) by additionally taking into account the energy of an electromagnetic field, making it the most general asymptotically flat and stationary solution of the Einstein–Maxwell equations in general relativity. As an electrovacuum solution, it only includes those charges associated with the magnetic field; it does not include any free electric charges.
Because observed astronomical objects do not possess an appreciable net electric charge (the magnetic fields of stars arise through other processes), the Kerr–Newman metric is primarily of theoretical interest.
The model lacks description of infalling baryonic matter, light (null dusts) or dark matter, and thus provides an incomplete description of stellar mass black holes and active galactic nuclei. The solution however is of mathematical interest and provides a fairly simple cornerstone for further exploration.
The Kerr–Newman solution is a special case of more general exact solutions of the Einstein–Maxwell equations with non-zero cosmological constant.
History
In December of 1963, Roy Kerr and Alfred Schild found the Kerr–Schild metrics that gave all Einstein spaces that are exact linear perturbations of Minkowski space. In early 1964, Kerr looked for all Einstein–Maxwell spaces with this same property. By February of 1964, the special case where the Kerr–Schild spaces were charged (including the Kerr–Newman solution) was known but the general case where the special directions were not geodesics of the underlying Minkowski space proved very difficult. The problem was given to George Debney to try to solve but was given up by March 1964. About this time Ezra T. Newman found the solution for charged Kerr by guesswork.
In 1965, Ezra "Ted" Newman found the axisymmetric solution of Einstein's field equation for a black hole which is both rotating and electrically charged. This formula for the metric tensor is called the Kerr–Newman metric. It is a generalisation of the Kerr metric for an uncharged spinning point-mass, which had been discovered by Roy Kerr two years earlier.
Four related solutions may be summarized by the following table:
{| class="wikitable"
|-
!
! Non-rotating (J = 0)
! Rotating (J ∈ )
|-
! Uncharged (Q = 0)
| Schwarzschild
| Kerr
|-
! Charged (Q ∈ )
| Reissner–Nordström
| Kerr-Newman
|-
|}
where Q represents the body's electric charge and J represents its spin angular momentum.
Overview of the solution
Newman's result represents the simplest stationary, axisymmetric, asymptotically flat solution of Einstein's equations in the presence of an electromagnetic field in four dimensions. It is sometimes referred to as an "electrovacuum" solution of Einstein's equations.
Any Kerr–Newman source has its rotation axis aligned with its magnetic axis. Thus, a Kerr–Newman source is different from commonly observed astronomical bodies, for which there is a substantial angle between the rotation axis and the magnetic moment. Specifically, neither the Sun, nor any of the planets in the Solar System have magnetic fields aligned with the spin axis. Thus, while the Kerr solution describes the gravitational field of the Sun and planets, the magnetic fields arise by a different process.
If the Kerr–Newman potential is considered as a model for a classical electron, it predicts an electron having not just a magnetic dipole moment, but also other multipole moments, such as an electric quadrupole moment. An electron quadrupole moment has not yet been experimentally detected; it appears to be zero.
In the G = 0 limit, the electromagnetic fields are those of a charged rotating disk inside a ring where the fields are infinite. The total field energy for this disk is infinite, and so this G = 0 limit does not solve the problem of infinite self-energy.
Like the Kerr metric for an uncharged rotating mass, the Kerr–Newman interior solution exists mathematically but is probably not representative of the actual metric of a physically realistic rotating black hole due to issues with the stability of the Cauchy horizon, due to mass inflation driven by infalling matter. Although it represents a generalization of the Kerr metric, it is not considered as very important for astrophysical purposes, since one does not expect that realistic black holes have a significant electric charge (they are expected to have a minuscule positive charge, but only because the proton has a much larger momentum than the electron, and is thus more likely to overcome electrostatic repulsion and be carried by momentum across the horizon).
The Kerr–Newman metric defines a black hole with an event horizon only when the combined charge and angular momentum are sufficiently small:
An electron's angular momentum J and charge Q (suitably specified in geometrized units) both exceed its mass M, in which case the metric has no event horizon. Thus, there can be no such thing as a black hole electron — only a naked spinning ring singularity. Such a metric has several seemingly unphysical properties, such as the ring's violation of the cosmic censorship hypothesis, and also appearance of causality-violating closed timelike curves in the immediate vicinity of the ring.
A 2009 paper by Russian theorist Alexander Burinskii considered an electron as a generalization of the previous models by Israel (1970) and Lopez (1984), which truncated the "negative" sheet of the Kerr-Newman metric, obtaining the source of the Kerr-Newman solution in the form of a relativistically rotating disk. Lopez's truncation regularized the Kerr-Newman metric by a cutoff at :, replacing the singularity by a flat regular space-time, the so called "bubble". Assuming that the Lopez bubble corresponds to a phase transition similar to the Higgs symmetry breaking mechanism, Burinskii showed that a gravity-created ring singularity forms by regularization the superconducting core of the electron model and should be described by the supersymmetric Landau-Ginzburg field model of phase transition:
By omitting Burinsky's intermediate work, we come to the recent new proposal: to consider the truncated by Israel and Lopez negative sheet of the KN solution as the sheet of the positron.
This modification unites the KN solution with the model of QED, and shows the important role of the Wilson lines formed by frame-dragging of the vector potential.
As a result, the modified KN solution acquires a strong interaction with Kerr's gravity caused by the additional energy contribution of the electron-positron vacuum and creates the Kerr–Newman relativistic circular string of Compton size.
Limiting cases
The Kerr–Newman metric can be seen to reduce to other exact solutions in general relativity in limiting cases. It reduces to
the Kerr metric as the charge Q goes to zero;
the Reissner–Nordström metric as the angular momentum J (or a = ) goes to zero;
the Schwarzschild metric as both the charge Q and the angular momentum J (or a) are taken to zero; and
Minkowski space if the mass M, the charge Q, and the rotational parameter a are all zero.
Alternately, if gravity is intended to be removed, Minkowski space arises if the gravitational constant G is zero, without taking the mass and charge to zero. In this case, the electric and magnetic fields are more complicated than simply the fields of a charged magnetic dipole; the zero-gravity limit is not trivial.
The metric
The Kerr–Newman metric describes the geometry of spacetime for a rotating charged black hole with mass M, charge Q and angular momentum J. The formula for this metric depends upon what coordinates or coordinate conditions are selected. Two forms are given below: Boyer–Lindquist coordinates, and Kerr–Schild coordinates. The gravitational metric alone is not sufficient to determine a solution to the Einstein field equations; the electromagnetic stress tensor must be given as well. Both are provided in each section.
Boyer–Lindquist coordinates
One way to express this metric is by writing down its line element in a particular set of spherical coordinates, also called Boyer–Lindquist coordinates:
where the coordinates are standard spherical coordinate system, and the length scales:
have been introduced for brevity. Here rs is the Schwarzschild radius of the massive body, which is related to its total mass-equivalent M by
where G is the gravitational constant, and rQ is a length scale corresponding to the electric charge Q of the mass
where ε0 is the vacuum permittivity.
Electromagnetic field tensor in Boyer–Lindquist form
The electromagnetic potential in Boyer–Lindquist coordinates is
while the Maxwell tensor is defined by
In combination with the Christoffel symbols the second order equations of motion can be derived with
where is the charge per mass of the test particle.
Kerr–Schild coordinates
The Kerr–Newman metric can be expressed in the Kerr–Schild form, using a particular set of Cartesian coordinates, proposed by Kerr and Schild in 1965. The metric is as follows.
Notice that k is a unit vector. Here M is the constant mass of the spinning object, Q is the constant charge of the spinning object, η is the Minkowski metric, and a = J/M is a constant rotational parameter of the spinning object. It is understood that the vector is directed along the positive z-axis, i.e. . The quantity r is not the radius, but rather is implicitly defined by the relation
Notice that the quantity r becomes the usual radius R
when the rotational parameter a approaches zero. In this form of solution, units are selected so that the speed of light is unity (c = 1). In order to provide a complete solution of the Einstein–Maxwell equations, the Kerr–Newman solution not only includes a formula for the metric tensor, but also a formula for the electromagnetic potential:
At large distances from the source (R ≫ a), these equations reduce to the Reissner–Nordström metric with:
In the Kerr–Schild form of the Kerr–Newman metric, the determinant of the metric tensor is everywhere equal to negative one, even near the source.
Electromagnetic fields in Kerr–Schild form
The electric and magnetic fields can be obtained in the usual way by differentiating the four-potential to obtain the electromagnetic field strength tensor. It will be convenient to switch over to three-dimensional vector notation.
The static electric and magnetic fields are derived from the vector potential and the scalar potential like this:
Using the Kerr–Newman formula for the four-potential in the Kerr–Schild form, in the limit of the mass going to zero, yields the following concise complex formula for the fields:
The quantity omega () in this last equation is similar to the Coulomb potential, except that the radius vector is shifted by an imaginary amount. This complex potential was discussed as early as the nineteenth century, by the French mathematician Paul Émile Appell.
Irreducible mass
The total mass-equivalent M, which contains the electric field-energy and the rotational energy, and the irreducible mass Mirr are related by
which can be inverted to obtain
In order to electrically charge and/or spin a neutral and static body, energy has to be applied to the system. Due to the mass–energy equivalence, this energy also has a mass-equivalent; therefore M is always higher than Mirr. If for example the rotational energy of a black hole is extracted via the Penrose processes, the remaining mass–energy will always stay greater than or equal to Mirr.
Important surfaces
Setting to 0 and solving for gives the inner and outer event horizon, which is located at the Boyer–Lindquist coordinate
Repeating this step with gives the inner and outer ergosphere
Equations of motion
For brevity, we further use nondimensionalized quantities normalized against , , and , where reduces to and to , and the equations of motion for a test particle of charge become
with for the total energy and for the axial angular momentum. is the Carter constant:
where is the poloidial component of the test particle's angular momentum, and the orbital inclination angle.
and
with and for particles are also conserved quantities.
is the frame dragging induced angular velocity. The shorthand term is defined by
The relation between the coordinate derivatives and the local 3-velocity is
for the radial,
for the poloidial,
for the axial and
for the total local velocity, where
is the axial radius of gyration (local circumference divided by 2π), and
the gravitational time dilation component. The local radial escape velocity for a neutral particle is therefore
| Physical sciences | Theory of relativity | Physics |
1448322 | https://en.wikipedia.org/wiki/Sodium%20bromide | Sodium bromide | Sodium bromide is an inorganic compound with the formula . It is a high-melting white, crystalline solid that resembles sodium chloride. It is a widely used source of the bromide ion and has many applications.
Synthesis, structure, reactions
NaBr crystallizes in the same cubic motif as NaCl, NaF and NaI. The anhydrous salt crystallizes above 50.7 °C. Dihydrate salt () crystallize out of water solution below 50.7 °C.
NaBr is produced by treating sodium hydroxide with hydrogen bromide.
Sodium bromide can be used as a source of the chemical element bromine. This can be accomplished by treating an aqueous solution of NaBr with chlorine gas:
Applications
Sodium bromide is the most useful inorganic bromide in industry. It is also used as a catalyst in TEMPO-mediated oxidation reactions.
Medicine
Also known as Sedoneural, sodium bromide has been used as a hypnotic, anticonvulsant, and sedative in medicine, widely used as an anticonvulsant and a sedative in the late 19th and early 20th centuries. Its action is due to the bromide ion, and for this reason potassium bromide is equally effective. In 1975, bromides were removed from drugs in the U.S. such as Bromo-Seltzer due to toxicity.
Preparation of other bromine compounds
Sodium bromide is widely used for the preparation of other bromides in organic synthesis and other areas. It is a source of the bromide nucleophile to convert alkyl chlorides to more reactive alkyl bromides by the Finkelstein reaction:
NaBr + RCl → RBr + NaCl (R = alkyl)
Once a large need in photography, but now shrinking, the photosensitive salt silver bromide is prepared using NaBr.
Disinfectant
Sodium bromide is used in conjunction with chlorine as a disinfectant for hot tubs and swimming pools.
Petroleum industry
Because of its high solubility in water (943.2 g/L or 9.16 mol/L, at 25 °C) sodium bromide is used to prepare dense drilling fluids used in oil wells to compensate a possible overpressure arising in the fluid column and to counteract the associated trend to blow out. The presence of the sodium cation also causes the bentonite added to the drilling fluid to swell, while the high ionic strength induces bentonite flocculation.
Safety
NaBr has a very low toxicity with an oral estimated at 3.5 g/kg for rats. However, this is a single-dose value. Bromide ion is a cumulative toxin with a relatively long half-life (in excess of a week in humans): see potassium bromide.
| Physical sciences | Halide salts | Chemistry |
1448472 | https://en.wikipedia.org/wiki/Laplacian%20matrix | Laplacian matrix | In the mathematical field of graph theory, the Laplacian matrix, also called the graph Laplacian, admittance matrix, Kirchhoff matrix or discrete Laplacian, is a matrix representation of a graph. Named after Pierre-Simon Laplace, the graph Laplacian matrix can be viewed as a matrix form of the negative discrete Laplace operator on a graph approximating the negative continuous Laplacian obtained by the finite difference method.
The Laplacian matrix relates to many useful properties of a graph. Together with Kirchhoff's theorem, it can be used to calculate the number of spanning trees for a given graph. The sparsest cut of a graph can be approximated through the Fiedler vector — the eigenvector corresponding to the second smallest eigenvalue of the graph Laplacian — as established by Cheeger's inequality. The spectral decomposition of the Laplacian matrix allows constructing low dimensional embeddings that appear in many machine learning applications and determines a spectral layout in graph drawing. Graph-based signal processing is based on the graph Fourier transform that extends the traditional discrete Fourier transform by substituting the standard basis of complex sinusoids for eigenvectors of the Laplacian matrix of a graph corresponding to the signal.
The Laplacian matrix is the easiest to define for a simple graph, but more common in applications for an edge-weighted graph, i.e., with weights on its edges — the entries of the graph adjacency matrix. Spectral graph theory relates properties of a graph to a spectrum, i.e., eigenvalues, and eigenvectors of matrices associated with the graph, such as its adjacency matrix or Laplacian matrix. Imbalanced weights may undesirably affect the matrix spectrum, leading to the need of normalization — a column/row scaling of the matrix entries — resulting in normalized adjacency and Laplacian matrices.
Definitions for simple graphs
Laplacian matrix
Given a simple graph with vertices , its Laplacian matrix is
defined element-wise as
or equivalently by the matrix
where D is the degree matrix and A is the adjacency matrix of the graph. Since is a simple graph, only contains 1s or 0s and its diagonal elements are all 0s.
Here is a simple example of a labelled, undirected graph and its Laplacian matrix.
We observe for the undirected graph that both the adjacency matrix and the Laplacian matrix are symmetric, and that row- and column-sums of the Laplacian matrix are all zeros (which directly implies that the Laplacian matrix is singular).
For directed graphs, either the indegree or outdegree might be used, depending on the application, as in the following example:
In the directed graph, both the adjacency matrix and the Laplacian matrix are asymmetric. In its Laplacian matrix, column-sums or row-sums are zero, depending on whether the indegree or outdegree has been used.
Laplacian matrix for an undirected graph via the oriented incidence matrix
The oriented incidence matrix B with element Bve for the vertex v and the edge e (connecting vertices and , with i ≠ j) is defined by
Even though the edges in this definition are technically directed, their directions can be arbitrary, still resulting in the same symmetric Laplacian matrix L defined as
where is the matrix transpose of B.
An alternative product defines the so-called edge-based Laplacian, as opposed to the original commonly used vertex-based Laplacian matrix L.
Symmetric Laplacian for a directed graph
The Laplacian matrix of a directed graph is by definition generally non-symmetric, while, e.g., traditional spectral clustering is primarily developed for undirected graphs with symmetric adjacency and Laplacian matrices. A trivial approach to apply techniques requiring the symmetry is to turn the original directed graph into an undirected graph and build the Laplacian matrix for the latter.
In the matrix notation, the adjacency matrix of the undirected graph could, e.g., be defined as a Boolean sum of the adjacency matrix of the original directed graph and its matrix transpose , where the zero and one entries of are treated as logical, rather than numerical, values, as in the following example:
Laplacian matrix normalization
A vertex with a large degree, also called a heavy node, results in a large diagonal entry in the Laplacian matrix dominating the matrix properties. Normalization is aimed to make the influence of such vertices more equal to that of other vertices, by dividing the entries of the Laplacian matrix by the vertex degrees. To avoid division by zero, isolated vertices with zero degrees are excluded from the process of the normalization.
Symmetrically normalized Laplacian
The symmetrically normalized Laplacian matrix is defined as:
where is the Moore–Penrose inverse.
The elements of are thus given by
The symmetrically normalized Laplacian matrix is symmetric if and only if the adjacency matrix is symmetric.
For a non-symmetric adjacency matrix of a directed graph, either of indegree and outdegree can be used for normalization:
Left (random-walk) and right normalized Laplacians
The left (random-walk) normalized Laplacian matrix is defined as:
where is the Moore–Penrose inverse.
The elements of are given by
Similarly, the right normalized Laplacian matrix is defined as
.
The left or right normalized Laplacian matrix is not symmetric if the adjacency matrix is symmetric, except for the trivial case of all isolated vertices. For example,
The example also demonstrates that if has no isolated vertices, then right stochastic and hence is the matrix of a random walk, so that the left normalized Laplacian has each row summing to zero. Thus we sometimes alternatively call the random-walk normalized Laplacian. In the less uncommonly used right normalized Laplacian each column sums to zero since is left stochastic.
For a non-symmetric adjacency matrix of a directed graph, one also needs to choose indegree or outdegree for normalization:
The left out-degree normalized Laplacian with row-sums all 0 relates to right stochastic , while the right in-degree normalized Laplacian with column-sums all 0 contains left stochastic .
Definitions for graphs with weighted edges
Common in applications graphs with weighted edges are conveniently defined by their adjacency matrices where values of the entries are numeric and no longer limited to zeros and ones. In spectral clustering and graph-based signal processing, where graph vertices represent data points, the edge weights can be computed, e.g., as inversely proportional to the distances between pairs of data points, leading to all weights being non-negative with larger values informally corresponding to more similar pairs of data points. Using correlation and anti-correlation between the data points naturally leads to both positive and negative weights. Most definitions for simple graphs are trivially extended to the standard case of non-negative weights, while negative weights require more attention, especially in normalization.
Laplacian matrix
The Laplacian matrix is defined by
where D is the degree matrix and A is the adjacency matrix of the graph.
For directed graphs, either the indegree or outdegree might be used, depending on the application, as in the following example:
Graph self-loops, manifesting themselves by non-zero entries on the main diagonal of the adjacency matrix, are allowed but do not affect the graph Laplacian values.
Symmetric Laplacian via the incidence matrix
For graphs with weighted edges one can define a weighted incidence matrix B and use it to construct the corresponding symmetric Laplacian as . An alternative cleaner approach, described here, is to separate the weights from the connectivity: continue using the incidence matrix as for regular graphs and introduce a matrix just holding the values of the weights. A spring system is an example of this model used in mechanics to describe a system of springs of given stiffnesses and unit length, where the values of the stiffnesses play the role of the weights of the graph edges.
We thus reuse the definition of the weightless incidence matrix B with element Bve for the vertex v and the edge e (connecting vertexes and , with i > j) defined by
We now also define a diagonal matrix W containing the edge weights. Even though the edges in the definition of B are technically directed, their directions can be arbitrary, still resulting in the same symmetric Laplacian matrix L defined as
where is the matrix transpose of B.
The construction is illustrated in the following example, where every edge is assigned the weight value i, with
Symmetric Laplacian for a directed graph
Just like for simple graphs, the Laplacian matrix of a directed weighted graph is by definition generally non-symmetric. The symmetry can be enforced by turning the original directed graph into an undirected graph first before constructing the Laplacian. The adjacency matrix of the undirected graph could, e.g., be defined as a sum of the adjacency matrix of the original directed graph and its matrix transpose as in the following example:
where the zero and one entries of are treated as numerical, rather than logical as for simple graphs, values, explaining the difference in the results - for simple graphs, the symmetrized graph still needs to be simple with its symmetrized adjacency matrix having only logical, not numerical values, e.g., the logical sum is 1 v 1 = 1, while the numeric sum is 1 + 1 = 2.
Alternatively, the symmetric Laplacian matrix can be calculated from the two Laplacians using the indegree and outdegree, as in the following example:
The sum of the out-degree Laplacian transposed and the in-degree Laplacian equals to the symmetric Laplacian matrix.
Laplacian matrix normalization
The goal of normalization is, like for simple graphs, to make the diagonal entries of the Laplacian matrix to be all unit, also scaling off-diagonal entries correspondingly. In a weighted graph, a vertex may have a large degree because of a small number of connected edges but with large weights just as well as due to a large number of connected edges with unit weights.
Graph self-loops, i.e., non-zero entries on the main diagonal of the adjacency matrix, do not affect the graph Laplacian values, but may need to be counted for calculation of the normalization factors.
Symmetrically normalized Laplacian
The symmetrically normalized Laplacian is defined as
where L is the unnormalized Laplacian, A is the adjacency matrix, D is the degree matrix, and is the Moore–Penrose inverse. Since the degree matrix D is diagonal, its reciprocal square root is just the diagonal matrix whose diagonal entries are the reciprocals of the square roots of the diagonal entries of D. If all the edge weights are nonnegative then all the degree values are automatically also nonnegative and so every degree value has a unique positive square root. To avoid the division by zero, vertices with zero degrees are excluded from the process of the normalization, as in the following example:
The symmetrically normalized Laplacian is a symmetric matrix if and only if the adjacency matrix A is symmetric and the diagonal entries of D are nonnegative, in which case we can use the term the symmetric normalized Laplacian.
The symmetric normalized Laplacian matrix can be also written as
using the weightless incidence matrix B and the diagonal matrix W containing the edge weights and defining the new weighted incidence matrix whose rows are indexed by the vertices and whose columns are indexed by the edges of G such that each column corresponding to an edge e = {u, v} has an entry in the row corresponding to u, an entry in the row corresponding to v, and has 0 entries elsewhere.
Random walk normalized Laplacian
The random walk normalized Laplacian is defined as
where D is the degree matrix. Since the degree matrix D is diagonal, its inverse is simply defined as a diagonal matrix, having diagonal entries which are the reciprocals of the corresponding diagonal entries of D. For the isolated vertices (those with degree 0), a common choice is to set the corresponding element to 0. The matrix elements of are given by
The name of the random-walk normalized Laplacian comes from the fact that this matrix is , where is simply the transition matrix of a random walker on the graph, assuming non-negative weights. For example, let denote the i-th standard basis vector. Then is a probability vector representing the distribution of a random walker's locations after taking a single step from vertex ; i.e., . More generally, if the vector is a probability distribution of the location of a random walker on the vertices of the graph, then is the probability distribution of the walker after steps.
The random walk normalized Laplacian can also be called the left normalized Laplacian since the normalization is performed by multiplying the Laplacian by the normalization matrix on the left. It has each row summing to zero since is right stochastic, assuming all the weights are non-negative.
In the less uncommonly used right normalized Laplacian each column sums to zero since is left stochastic.
For a non-symmetric adjacency matrix of a directed graph, one also needs to choose indegree or outdegree for normalization:
The left out-degree normalized Laplacian with row-sums all 0 relates to right stochastic , while the right in-degree normalized Laplacian with column-sums all 0 contains left stochastic .
Negative weights
Negative weights present several challenges for normalization:
The presence of negative weights may naturally result in zero row- and/or column-sums for non-isolated vertices. A vertex with a large row-sum of positive weights and equally negatively large row-sum of negative weights, together summing up to zero, could be considered a heavy node and both large values scaled, while the diagonal entry remains zero, like for an isolated vertex.
Negative weights may also give negative row- and/or column-sums, so that the corresponding diagonal entry in the non-normalized Laplacian matrix would be negative and a positive square root needed for the symmetric normalization would not exist.
Arguments can be made to take the absolute value of the row- and/or column-sums for the purpose of normalization, thus treating a possible value -1 as a legitimate unit entry of the main diagonal of the normalized Laplacian matrix.
Properties
For an (undirected) graph G and its Laplacian matrix L with eigenvalues :
L is symmetric.
L is positive-semidefinite (that is for all ). This can be seen from the fact that the Laplacian is symmetric and diagonally dominant.
L is an M-matrix (its off-diagonal entries are nonpositive, yet the real parts of its eigenvalues are nonnegative).
Every row sum and column sum of L is zero. Indeed, in the sum, the degree of the vertex is summed with a "−1" for each neighbor.
In consequence, , because the vector satisfies This also implies that the Laplacian matrix is singular.
The number of connected components in the graph is the dimension of the nullspace of the Laplacian and the algebraic multiplicity of the 0 eigenvalue.
The smallest non-zero eigenvalue of L is called the spectral gap.
The second smallest eigenvalue of L (could be zero) is the algebraic connectivity (or Fiedler value) of G and approximates the sparsest cut of a graph.
The Laplacian is an operator on the n-dimensional vector space of functions , where is the vertex set of G, and .
When G is k-regular, the normalized Laplacian is: , where A is the adjacency matrix and I is an identity matrix.
For a graph with multiple connected components, L is a block diagonal matrix, where each block is the respective Laplacian matrix for each component, possibly after reordering the vertices (i.e. L is permutation-similar to a block diagonal matrix).
The trace of the Laplacian matrix L is equal to where is the number of edges of the considered graph.
Now consider an eigendecomposition of , with unit-norm eigenvectors and corresponding eigenvalues :
Because can be written as the inner product of the vector with itself, this shows that and so the eigenvalues of are all non-negative.
All eigenvalues of the normalized symmetric Laplacian satisfy 0 = μ0 ≤ … ≤ μn−1 ≤ 2. These eigenvalues (known as the spectrum of the normalized Laplacian) relate well to other graph invariants for general graphs.
One can check that:
,
i.e., is similar to the normalized Laplacian . For this reason, even if is in general not symmetric, it has real eigenvalues — exactly the same as the eigenvalues of the normalized symmetric Laplacian .
Interpretation as the discrete Laplace operator approximating the continuous Laplacian
The graph Laplacian matrix can be further viewed as a matrix form of the negative discrete Laplace operator on a graph approximating the negative continuous Laplacian operator obtained by the finite difference method.
(See Discrete Poisson equation) In this interpretation, every graph vertex is treated as a grid point; the local connectivity of the vertex determines the finite difference approximation stencil at this grid point, the grid size is always one for every edge, and there are no constraints on any grid points, which corresponds to the case of the homogeneous Neumann boundary condition, i.e., free boundary. Such an interpretation allows one, e.g., generalizing the Laplacian matrix to the case of graphs with an infinite number of vertices and edges, leading to a Laplacian matrix of an infinite size.
Generalizations and extensions of the Laplacian matrix
Generalized Laplacian
The generalized Laplacian is defined as:
Notice the ordinary Laplacian is a generalized Laplacian.
Admittance matrix of an AC circuit
The Laplacian of a graph was first introduced to model electrical networks.
In an alternating current (AC) electrical network, real-valued resistances are replaced by complex-valued impedances.
The weight of edge (i, j) is, by convention, minus the reciprocal of the impedance directly between i and j.
In models of such networks, the entries of the adjacency matrix are complex, but the Kirchhoff matrix remains symmetric, rather than being Hermitian.
Such a matrix is usually called an "admittance matrix", denoted , rather than a "Laplacian".
This is one of the rare applications that give rise to complex symmetric matrices.
Magnetic Laplacian
There are other situations in which entries of the adjacency matrix are complex-valued, and the Laplacian does become a Hermitian matrix. The Magnetic Laplacian for a directed graph with real weights is constructed as the Hadamard product of the real symmetric matrix of the symmetrized Laplacian and the Hermitian phase matrix with the complex entries
which encode the edge direction into the phase in the complex plane.
In the context of quantum physics, the magnetic Laplacian can be interpreted as the operator that describes the phenomenology of a free charged particle on a graph, which is subject to the action of a magnetic field and the parameter is called electric charge.
In the following example :
Deformed Laplacian
The deformed Laplacian is commonly defined as
where I is the identity matrix, A is the adjacency matrix, D is the degree matrix, and s is a (complex-valued) number. The standard Laplacian is just and is the signless Laplacian.
Signless Laplacian
The signless Laplacian is defined as
where is the degree matrix, and is the adjacency matrix. Like the signed Laplacian , the signless Laplacian also is positive semi-definite as it can be factored as
where is the incidence matrix. has a 0-eigenvector if and only if it has a bipartite connected component (isolated vertices being bipartite connected components). This can be shown as
This has a solution where if and only if the graph has a bipartite connected component.
Directed multigraphs
An analogue of the Laplacian matrix can be defined for directed multigraphs. In this case the Laplacian matrix L is defined as
where D is a diagonal matrix with Di,i equal to the outdegree of vertex i and A is a matrix with Ai,j equal to the number of edges from i to j (including loops).
Open source software implementations
SciPy
NetworkX
Julia
Application software
scikit-learn Spectral Clustering
PyGSP: Graph Signal Processing in Python
megaman: Manifold Learning for Millions of Points
smoothG
Laplacian Change Point Detection for Dynamic Graphs (KDD 2020)
LaplacianOpt (A Julia Package for Maximizing Laplacian's Second Eigenvalue of Weighted Graphs)
LigMG (Large Irregular Graph MultiGrid)
Laplacians.jl
| Mathematics | Graph theory | null |
1449031 | https://en.wikipedia.org/wiki/Activity%20coefficient | Activity coefficient | In thermodynamics, an activity coefficient is a factor used to account for deviation of a mixture of chemical substances from ideal behaviour. In an ideal mixture, the microscopic interactions between each pair of chemical species are the same (or macroscopically equivalent, the enthalpy change of solution and volume variation in mixing is zero) and, as a result, properties of the mixtures can be expressed directly in terms of simple concentrations or partial pressures of the substances present e.g. Raoult's law. Deviations from ideality are accommodated by modifying the concentration by an activity coefficient. Analogously, expressions involving gases can be adjusted for non-ideality by scaling partial pressures by a fugacity coefficient.
The concept of activity coefficient is closely linked to that of activity in chemistry.
Thermodynamic definition
The chemical potential, , of a substance B in an ideal mixture of liquids or an ideal solution is given by
,
where μ is the chemical potential of a pure substance , and is the mole fraction of the substance in the mixture.
This is generalised to include non-ideal behavior by writing
when is the activity of the substance in the mixture,
,
where is the activity coefficient, which may itself depend on . As approaches 1, the substance behaves as if it were ideal. For instance, if ≈ 1, then Raoult's law is accurate. For > 1 and < 1, substance B shows positive and negative deviation from Raoult's law, respectively. A positive deviation implies that substance B is more volatile.
In many cases, as goes to zero, the activity coefficient of substance B approaches a constant; this relationship is Henry's law for the solvent. These relationships are related to each other through the Gibbs–Duhem equation.
Note that in general activity coefficients are dimensionless.
In detail: Raoult's law states that the partial pressure of component B is related to its vapor pressure (saturation pressure) and its mole fraction in the liquid phase,
with the convention
In other words: Pure liquids represent the ideal case.
At infinite dilution, the activity coefficient approaches its limiting value, ∞. Comparison with Henry's law,
immediately gives
In other words: The compound shows nonideal behavior in the dilute case.
The above definition of the activity coefficient is impractical if the compound does not exist as a pure liquid. This is often the case for electrolytes or biochemical compounds. In such cases, a different definition is used that considers infinite dilution as the ideal state:
with
and
The symbol has been used here to distinguish between the two kinds of activity coefficients. Usually it is omitted, as it is clear from the context which kind is meant. But there are cases where both kinds of activity coefficients are needed and may even appear in the same equation, e.g., for solutions of salts in (water + alcohol) mixtures. This is sometimes a source of errors.
Modifying mole fractions or concentrations by activity coefficients gives the effective activities of the components, and hence allows expressions such as Raoult's law and equilibrium constants to be applied to both ideal and non-ideal mixtures.
Knowledge of activity coefficients is particularly important in the context of electrochemistry since the behaviour of electrolyte solutions is often far from ideal, due to the effects of the ionic atmosphere. Additionally, they are particularly important in the context of soil chemistry due to the low volumes of solvent and, consequently, the high concentration of electrolytes.
Ionic solutions
For solution of substances which ionize in solution the activity coefficients of the cation and anion cannot be experimentally determined independently of each other because solution properties depend on both ions. Single ion activity coefficients must be linked to the activity coefficient of the dissolved electrolyte as if undissociated. In this case a mean stoichiometric activity coefficient of the dissolved electrolyte, γ±, is used. It is called stoichiometric because it expresses both the deviation from the ideality of the solution and the incomplete ionic dissociation of the ionic compound which occurs especially with the increase of its concentration.
For a 1:1 electrolyte, such as NaCl it is given by the following:
where and are the activity coefficients of the cation and anion respectively.
More generally, the mean activity coefficient of a compound of formula is given by
Single-ion activity coefficients can be calculated theoretically, for example by using the Debye–Hückel equation. The theoretical equation can be tested by combining the calculated single-ion activity coefficients to give mean values which can be compared to experimental values.
The prevailing view that single ion activity coefficients are unmeasurable independently, or perhaps even physically meaningless, has its roots in the work of Guggenheim in the late 1920s. However, chemists have never been able to give up the idea of single ion activities, and by implication single ion activity coefficients. For example, pH is defined as the negative logarithm of the hydrogen ion activity. If the prevailing view on the physical meaning and measurability of single ion activities is correct then defining pH as the negative logarithm of the hydrogen ion activity places the quantity squarely in the unmeasurable category. Recognizing this logical difficulty, International Union of Pure and Applied Chemistry (IUPAC) states that the activity-based definition of pH is a notional definition only. Despite the prevailing negative view on the measurability of single ion coefficients, the concept of single ion activities continues to be discussed in the literature.
Concentrated ionic solutions
For concentrated ionic solutions the hydration of ions must be taken into consideration, as done by Stokes and Robinson in their hydration model from 1948. The activity coefficient of the electrolyte is split into electric and statistical components by E. Glueckauf who modifies the Robinson–Stokes model.
The statistical part includes hydration index number , the number of ions from the dissociation and the ratio between the apparent molar volume of the electrolyte and the molar volume of water and molality .
Concentrated solution statistical part of the activity coefficient is:
The Stokes–Robinson model has been analyzed and improved by other investigators. The problem with this widely accepted idea that electrolyte activity coefficients are driven at higher concentrations by changes in hydration is that water activities are completely dependent on the concentration of the ions themselves, as imposed by a thermodynamic relationship called the Gibbs-Duhem equation. This means that the activity coefficients and the corresponding water activities are linked together fundamentally, regardless of molecular-level hypotheses. Due to this high correlation, such hypotheses are not independent enough to be satisfactorily tested.
The rise in activity coefficients found with most aqueous strong electrolyte systems can be explained more plausibly by increasing electrostatic repulsions between ions of the same charge which are forced together as the available space between them decreases. In this way, the initial attractions between cations and anions at the low concentrations described by Debye and Hueckel are progressively overcome. It has been proposed that these electrostatic repulsions take place predominantly through the formation of so-called ion trios in which two ions of like charge interact, on average and at distance, with the same counterion as well as with each other. This model accurately reproduces the experimental patterns of activity and osmotic coefficients exhibited by numerous 3-ion aqueous electrolyte mixtures.
Experimental determination of activity coefficients
Activity coefficients may be determined experimentally by making measurements on non-ideal mixtures. Use may be made of Raoult's law or Henry's law to provide a value for an ideal mixture against which the experimental value may be compared to obtain the activity coefficient. Other colligative properties, such as osmotic pressure may also be used.
Radiochemical methods
Activity coefficients can be determined by radiochemical methods.
At infinite dilution
Activity coefficients for binary mixtures are often reported at the infinite dilution of each component. Because activity coefficient models simplify at infinite dilution, such empirical values can be used to estimate interaction energies. Examples are given for water:
Theoretical calculation of activity coefficients
Activity coefficients of electrolyte solutions may be calculated theoretically, using the Debye–Hückel equation or extensions such as the Davies equation, Pitzer equations or TCPC model. Specific ion interaction theory (SIT) may also be used.
For non-electrolyte solutions correlative methods such as UNIQUAC, NRTL, MOSCED or UNIFAC may be employed, provided fitted component-specific or model parameters are available. COSMO-RS is a theoretical method which is less dependent on model parameters as required information is obtained from quantum mechanics calculations specific to each molecule (sigma profiles) combined with a statistical thermodynamics treatment of surface segments.
For uncharged species, the activity coefficient γ0 mostly follows a salting-out model:
This simple model predicts activities of many species (dissolved undissociated gases such as CO2, H2S, NH3, undissociated acids and bases) to high ionic strengths (up to 5 mol/kg). The value of the constant b for CO2 is 0.11 at 10 °C and 0.20 at 330 °C.
For water as solvent, the activity aw can be calculated using:
where ν is the number of ions produced from the dissociation of one molecule of the dissolved salt, b is the molality of the salt dissolved in water, φ is the osmotic coefficient of water, and the constant 55.51 represents the molality of water. In the above equation, the activity of a solvent (here water) is represented as inversely proportional to the number of particles of salt versus that of the solvent.
Link to ionic diameter
The ionic activity coefficient is connected to the ionic diameter by the formula obtained from Debye–Hückel theory of electrolytes:
where A and B are constants, zi is the valence number of the ion, and I is ionic strength.
Dependence on state parameters
The derivative of an activity coefficient with respect to temperature is related to excess molar enthalpy by
Similarly, the derivative of an activity coefficient with respect to pressure can be related to excess molar volume.
Application to chemical equilibrium
At equilibrium, the sum of the chemical potentials of the reactants is equal to the sum of the chemical potentials of the products. The Gibbs free energy change for the reactions, ΔrG, is equal to the difference between these sums and therefore, at equilibrium, is equal to zero. Thus, for an equilibrium such as
Substitute in the expressions for the chemical potential of each reactant:
Upon rearrangement this expression becomes
The sum
is the standard free energy change for the reaction, .
Therefore,
where is the equilibrium constant. Note that activities and equilibrium constants are dimensionless numbers.
This derivation serves two purposes. It shows the relationship between standard free energy change and equilibrium constant. It also shows that an equilibrium constant is defined as a quotient of activities. In practical terms this is inconvenient. When each activity is replaced by the product of a concentration and an activity coefficient, the equilibrium constant is defined as
where [S] denotes the concentration of S, etc. In practice equilibrium constants are determined in a medium such that the quotient of activity coefficients is constant and can be ignored, leading to the usual expression
which applies under the conditions that the activity quotient has a particular (constant) value.
| Physical sciences | Thermodynamics | Chemistry |
1449175 | https://en.wikipedia.org/wiki/Degree%20matrix | Degree matrix | In the mathematical field of algebraic graph theory, the degree matrix of an undirected graph is a diagonal matrix which contains information about the degree of each vertex—that is, the number of edges attached to each vertex. It is used together with the adjacency matrix to construct the Laplacian matrix of a graph: the Laplacian matrix is the difference of the degree matrix and the adjacency matrix.
Definition
Given a graph with , the degree matrix for is a diagonal matrix defined as
where the degree of a vertex counts the number of times an edge terminates at that vertex. In an undirected graph, this means that each loop increases the degree of a vertex by two. In a directed graph, the term degree may refer either to indegree (the number of incoming edges at each vertex) or outdegree (the number of outgoing edges at each vertex).
Example
The following undirected graph has a 6x6 degree matrix with values:
Note that in the case of undirected graphs, an edge that starts and ends in the same node increases the corresponding degree value by 2 (i.e. it is counted twice).
Properties
The degree matrix of a k-regular graph has a constant diagonal of .
According to the degree sum formula, the trace of the degree matrix is twice the number of edges of the considered graph.
| Mathematics | Graph theory | null |
1449823 | https://en.wikipedia.org/wiki/Plotopteridae | Plotopteridae | Plotopteridae is an extinct family of flightless seabirds with uncertain placement, generally considered as member of order Suliformes. They exhibited remarkable convergent evolution with the penguins, particularly with the now extinct giant penguins. That they lived in the North Pacific, the other side of the world from the penguins, has led to them being described at times as the Northern Hemisphere's penguins, though they were not closely related. More recent studies have shown, however, that the shoulder-girdle, forelimb and sternum of plotopterids differ significantly from those of penguins, so comparisons in terms of function may not be entirely accurate. Plotopterids are regarded as closely related to Anhingidae (darters) and Phalacrocoracidae (cormorants). On the other hand, there is a theory that this group may have a common ancestor with penguins due to the similarity of forelimb and brain morphology. However, the endocast morphology of stem group Sphenisciformes differs from both Plotopteridae and modern penguins.
Their fossils have been found in California, Oregon, Washington, British Columbia, Hokkaido, Tōhoku, Chūbu, Kyushu. They seem to have evolved on arctic islands during the mid-Eocene, spreading southwards with the formation of kelp forests They ranged in size from that of a large cormorant (such as a Brandt's cormorant), to very large size, with femur length two times longer than emperor penguin. They had shortened wings optimised for underwater wing-propelled pursuit diving (like penguins or the now extinct great auk), and a body skeleton similar to that of the darter.
The second species to be named from rocks along the eastern Pacific Ocean was Tonsala hildegardae from the late Oligocene lower part of the Pysht Formation in Washington State. More fossils of T. hildegardae have since been described and included some of the first known examples of borings made by the marine bone-eating worm Osedax in bird bones.
The earliest known member of the family, Phocavis maritimus lived in the late Eocene, but most of the known species lived during Oligocene time, becoming extinct in the early to mid-Miocene. That they became extinct at the same time as the giant penguins of the Southern Hemisphere, which also coincided with the radiation of the seals and dolphins, has led to speculation that the expansion of marine mammals was responsible for the extinction of the Plotopteridae, though this has not been formally tested.
| Biology and health sciences | Prehistoric birds | Animals |
1449983 | https://en.wikipedia.org/wiki/Natural%20reservoir | Natural reservoir | In infectious disease ecology and epidemiology, a natural reservoir, also known as a disease reservoir or a reservoir of infection, is the population of organisms or the specific environment in which an infectious pathogen naturally lives and reproduces, or upon which the pathogen primarily depends for its survival. A reservoir is usually a living host of a certain species, such as an animal or a plant, inside of which a pathogen survives, often (though not always) without causing disease for the reservoir itself. By some definitions a reservoir may also be an environment external to an organism, such as a volume of contaminated air or water.
Because of the enormous variety of infectious microorganisms capable of causing disease, precise definitions for what constitutes a natural reservoir are numerous, various, and often conflicting. The reservoir concept applies only for pathogens capable of infecting more than one host population and only with respect to a defined target population – the population of organisms in which the pathogen causes disease. The reservoir is any population of organisms (or any environment) which harbors the pathogen and transmits it to the target population. Reservoirs may comprise one or more different species, may be the same or a different species as the target, and, in the broadest sense, may include vector species, which are otherwise distinct from natural reservoirs. Significantly, species considered reservoirs for a given pathogen may not experience symptoms of disease when infected by the pathogen.
Identifying the natural reservoirs of infectious pathogens has proven useful in treating and preventing large outbreaks of disease in humans and domestic animals, especially those diseases for which no vaccine exists. In principle, zoonotic diseases can be controlled by isolating or destroying the pathogen's reservoirs of infection. The mass culling of animals confirmed or suspected as reservoirs for human pathogens, such as birds that harbor avian influenza, has been effective at containing possible epidemics in many parts of the world; for other pathogens, such as the ebolaviruses, the identity of the presumed natural reservoir remains obscure.
Definition and terminology
The great diversity of infectious pathogens, their possible hosts, and the ways in which their hosts respond to infection has resulted in multiple definitions for "natural reservoir", many of which are conflicting or incomplete. In a 2002 conceptual exploration published in the CDC's Emerging Infectious Diseases, the natural reservoir of a given pathogen is defined as "one or more epidemiologically connected populations or environments in which the pathogen can be permanently maintained and from which infection is transmitted to the defined target population." The target population is the population or species in which the pathogen causes disease; it is the population of interest because it has disease when infected by the pathogen (for example, humans are the target population in most medical epidemiological studies).
A common criterion in other definitions distinguishes reservoirs from non-reservoirs by the degree to which the infected host shows symptoms of disease. By these definitions, a reservoir is a host that does not experience the symptoms of disease when infected by the pathogen, whereas non-reservoirs show symptoms of the disease. The pathogen still feeds, grows, and reproduces inside a reservoir host, but otherwise does not significantly affect its health; the relationship between pathogen and reservoir is more or less commensal, whereas in susceptible hosts that do develop disease caused by the pathogen, the pathogen is considered parasitic.
What further defines a reservoir for a specific pathogen is where it can be maintained and from where it can be transmitted. A "multi-host" organism is capable of having more than one natural reservoir.
Types of reservoirs
Natural reservoirs can be divided into three main types: human, animal (non-human), and environmental.
Human reservoirs
Human reservoirs are human beings infected by pathogens that exist on or within the human body. Infections like poliomyelitis and smallpox, which exist exclusively within a human reservoir, are sometimes known as anthroponoses. Humans can act as reservoirs for sexually transmitted diseases, measles, mumps, streptococcal infection, various respiratory pathogens, and the smallpox virus.
Animal reservoirs
Animal (non-human) reservoirs consist of domesticated and wild animals infected by pathogens. For example, the bacterium Vibrio cholerae, which causes cholera in humans, has natural reservoirs in copepods, zooplankton, and shellfish. Parasitic blood-flukes of the genus Schistosoma, responsible for schistosomiasis, spend part of their lives inside freshwater snails before completing their life cycles in vertebrate hosts. Viruses of the taxon Ebolavirus, which causes Ebola virus disease, are thought to have a natural reservoir in bats or other animals exposed to the virus. Other zoonotic diseases that have been transmitted from animals to humans include: rabies, blastomycosis, psittacosis, trichinosis, cat-scratch disease, histoplasmosis, coccidioidomycosis, and salmonella.
Common animal reservoirs include: bats, rodents, cows, pigs, sheep, swine, rabbits, raccoons, dogs, and other mammals.
Common animal reservoirs
Bats
Numerous zoonotic diseases have been traced back to bats. There are a couple of theories that serve as possible explanations as to why bats carry so many viruses. One proposed theory is that there exist so many bat-borne illnesses because there exist a large number of bat species and individuals. The second possibility is that something about bats' physiology makes them especially good reservoir hosts. Perhaps bats' "food choices, population structure, ability to fly, seasonal migration and daily movement patterns, torpor and hibernation, life span, and roosting behaviors" are responsible for making them especially suitable reservoir hosts. Lyssaviruses (including the Rabies virus), Henipaviruses, Menangle and Tioman viruses, SARS-CoV-Like Viruses, and Ebola viruses have all been traced back to different species of bats. Fruit bats in particular serve as the reservoir host for Nipah virus (NiV).
Rats
Rats are known to be the reservoir hosts for several zoonotic diseases. Norway rats were found to be infested with the Lyme disease spirochetes. In Mexico rats are known carriers of Trypanosoma cruzi, which causes Chagas disease.
Mice
White-footed mice (Peromyscus leucopus) are one of the most important animal reservoirs for the Lyme disease spirochete (Borrelia burgdorferi). Deer mice serve as reservoir hosts for Sin Nombre virus, which causes hantavirus pulmonary syndrome (HPS).
Monkeys
The Zika virus originated from monkeys in Africa. In São José do Rio Preto and Belo Horizonte, Brazil the Zika virus has been found in dead monkeys. Genome sequencing has revealed the virus to be very similar to the type that infects humans.
Environmental reservoirs
Environmental reservoirs include living and non-living reservoirs that harbor infectious pathogens outside the bodies of animals. These reservoirs may exist on land (plants and soil), in water, or the air. Pathogens in these reservoirs are sometimes free-living. The bacteria Legionella pneumophila, a facultative intracellular parasite which causes Legionnaires' disease, and Vibrio cholerae, which causes cholera, can both exist as free-living parasites in certain water sources as well as in invertebrate animal hosts.
Disease transmission
A disease reservoir acts as a transmission point between a pathogen and a susceptible host. Transmission can occur directly or indirectly.
Direct transmission
Direct transmission can occur from direct contact or direct droplet spread. Direct contact transmission between two people can happen through skin contact, kissing, and sexual contact. Humans serving as disease reservoirs can be symptomatic (showing illness) or asymptomatic (not showing illness), act as disease carriers, and often spread illness unknowingly. Human carriers commonly transmit disease because they do not realize they are infected and take no special precautions to prevent transmission. Symptomatic persons aware of their illness are not as likely to transmit infection because they take precautions to reduce possible transmission of the disease and/or seek out treatment to prevent the spread of the disease. Direct droplet spread is due to solid particles or liquid droplets suspended in the air for some time. Droplet spread is considered the transmission of the pathogen to a susceptible host within a meter of distance; said droplet spread can occur from coughing, sneezing, and/or just talking.
Neisseria gonorrhoeae (Gonorrhea) is transmitted by sexual contact involving the penis, vagina, mouth, and anus through direct contact transmission.
Bordetella pertussis (Pertussis) is transmitted by cough from the human reservoir to the susceptible host through direct droplet spread.
Indirect transmission
Indirect transmission can occur by airborne transmission, by vehicles (including fomites), and by vectors.
Airborne transmission is different from direct droplet spread as it is defined as disease transmission that takes place over a distance larger than a meter. Pathogens that can be transmitted through airborne sources are carried by particles such as dust or dried residue (referred to as droplet nuclei).
Vehicles such as food, water, blood and fomites can act as passive transmission points between reservoirs and susceptible hosts. Fomites are inanimate objects (doorknobs, medical equipment, etc.) that become contaminated by a reservoir source or someone/something that is a carrier. A vehicle, like a reservoir, may also be a favorable environment for the growth of an infectious agent, as coming into contact with a vehicle leads to its transmission.
Vector transmission occurs most often from insect bites from mosquitoes, flies, fleas, and ticks. There are two sub-categories of vectors: mechanical (an insect transmits the pathogen to a host without the insect itself being affected) and biological (reproduction of the pathogen occurs within the vector before the pathogen is transmitted to a host). To give a few examples, Morbillivirus (measles) is transmitted from an infected human host to a susceptible host as they are transmitted by respiration through airborne transmission. Campylobacter (campylobacteriosis) is a common bacterial infection that is spread from human or non-human reservoirs by vehicles such as contaminated food and water. Plasmodium falciparum (malaria) can be transmitted from an infected mosquito, an animal (non-human) reservoir, to a human host by biological vector transmission.
Implications for public health
LH Taylor found that 61% of all human pathogens are classified as zoonotic. Thus, the identification of the natural reservoirs of pathogens before zoonosis would be incredibly useful from a public health standpoint. Preventive measures can be taken to lessen the frequency of outbreaks, such as vaccinating the animal sources of disease or preventing contact with reservoir host animals. To predict and prevent future outbreaks of zoonotic diseases, the U.S. Agency for International Development started the Emerging Pandemic Threats initiative in 2009. In alliance with University of California-Davis, EcoHealth Alliance, Metabiota Inc., Smithsonian Institution, and Wildlife Conservation Society with support from Columbia and Harvard universities, the members of the PREDICT project are focusing on the "detection and discovery of zoonotic diseases at the wildlife-human interface." There are numerous other organizations around the world experimenting with different methods to predict and identify reservoir hosts. Researchers at the University of Glasgow created a machine learning algorithm that is designed to use "viral genome sequences to predict the likely natural host for a broad spectrum of RNA viruses, the viral group that most often jumps from animals to humans."
| Biology and health sciences | Concepts | Health |
1450081 | https://en.wikipedia.org/wiki/Sequence%20stratigraphy | Sequence stratigraphy | Sequence stratigraphy is a branch of geology, specifically a branch of stratigraphy, that attempts to discern and understand historic geology through time by subdividing and linking sedimentary deposits into unconformity bounded units on a variety of scales. The essence of the method is mapping of strata based on identification of surfaces which are assumed to represent time lines (e.g. subaerial unconformities, maximum flooding surfaces), thereby placing stratigraphy in chronostratigraphic framework allowing understanding of the evolution of the Earth's surface in a particular region through time. Sequence stratigraphy is a useful alternative to a purely lithostratigraphic approach, which emphasizes solely based on the compositional similarity of the lithology of rock units rather than time significance. Unconformities are particularly important in understanding geologic history because they represent erosional surfaces where there is a clear gap in the record. Conversely within a sequence the geologic record should be relatively continuous and complete record that is genetically related.
Stratigraphers explain sequence boundaries and stratigraphic units primarily in terms of changes in relative sea level (the combination of global changes in eustatic sea level and regional subsidence caused by tectonic subsidence, thermal subsidence and load-induced subsidence as the weight of accumulated sediment and water cause isostatic subsidence as a sedimentary basin is filled). The net changes resulting from these vertical forces increases or reduces accommodation space for sediments to accumulate in a sedimentary basin. A secondary influence is the rate of sediment supply to the basin which determines the rate at which that space is filled.
Historical development
The origin of sequence stratigraphy can be traced back to the work of L.L. Sloss on interregional unconformities of the North American craton. Sloss recognized six craton-wide sequences representing hundreds of millions of years of earth history. In the late 1960s Sloss had several students, notably Peter Vail, Robert Mitchum, and John Sangree, who completed dissertations studying the Pennsylvanian sedimentary rocks of the North American craton and became aware that global changes in sea level could have been responsible for the numerous widespread unconformities in those rocks. During their subsequent careers as research scientists at Exxon's research division Vail, Mitchum and others pioneered the practice of seismic stratigraphy, the stratigraphic interpretation of seismic reflection profiles to understand the layering and packaging of sedimentary rocks in the subsurface using acoustic imaging. The advent of seismic stratigraphy made it possible to identify sequences representing shorter period of time ranging in duration from tens of thousands to a few million years; and to compare the sequence stratigraphic history around the globe. This in turn led to sequence stratigraphy becoming systematized and understood to have widespread application to stratigraphic study of rock outcrops on the earth's surface as well. During the 1980s this ushered in a revolution in stratigraphy based on the delineation of regional physical surfaces that separate the sedimentary rock into packages representing discrete and sequential periods of time and predictable patterns of sediment depositional history.
Significant surfaces
Sequence boundaries
Sequence boundaries are deemed the most significant surfaces. Sequence boundaries are defined as unconformities or their correlative conformities. Sequence boundaries are formed due to the sea level fall. For example, multi-story fluvial sandstone packages often infill incised valleys formed by the sea level drop associated with sequence boundaries. The incised valleys of sequence boundaries correlate laterally with interfluves, palaeosols formed on the margins of incised valleys. The valley infills are not genetically related to underlying depositional systems as previous interpretations thought. There are four criteria distinguishing incised valley fills from other types of multi-story sandstone deposits: a widespread correlation with a regional, high relief erosional surface that is more widespread than the erosional bases of individual channels within the valley; facies associations reflect a basinward shift in facies when compared with underlying units; erosional base of the valley removes preceding systems tracts and marine bands producing a time gap, the removed units will be preserved beneath the interfluves; increasing channel fill and fine grained units upwards or changes in the character of the fluvial systems reflecting increasing accommodation space. Sandstone bodies associated with incised valleys can be good hydrocarbon reservoirs. There have been problems in the correlation and distribution of these bodies. Sequence stratigraphic principles and identification of significant surfaces have resolved some issues.
Parasequence boundaries
Lesser importance is attached to parasequence boundaries, however, there is a suggestion that flooding surfaces representing parasequence boundaries may be more laterally extensive leaving more evidence than sequence boundaries because the coastal plain has a lower gradient than the inner continental shelf. Parasequence boundaries may be distinguished by differences in physical and chemical properties across the surface such as; formation water salinity, hydrocarbon properties, porosity, compressional velocities and mineralogy. Parasequence boundaries may not form a barrier to hydrocarbon accumulation but may inhibit vertical reservoir communication. After production begins the parasequences act as separate drainage units with the flooding surfaces, which are overlain by shales or carbonate-cemented horizons, forming a barrier to vertical reservoir communications. Sequence stratigraphic principles have optimized production potential once reservoir scale architecture is identified and separate drainage units identified.
Systems tracts
The concept of systems tracts evolved to link the contemporaneous depositional systems. Systems tracts form subdivision in a sequence. Different kinds of systems tracts are assigned on the basis of stratal stacking pattern, position in a sequence, and in the sea level curve and types of bounding surfaces.
A lowstand systems tract (LST) forms when the rate of sedimentation outpaces the rate of sea level rise during the early stage of the sea level curve. It is bounded by a subaerial unconformity or its correlative conformity at the base and maximum regressive surface at the top.
A transgressive systems tract (TST) is bounded by maximum regressive surface at the base and maximum flooding surface at the top. This systems tracts forms when the rate of sedimentation is outpaced by the rate of sea level rise in the sea level curves.
A highstand systems tract (HST) occurs during the late stage of base level rise when the rate of sea level rise drops below the sedimentation rate. In this period of sea level highstand is formed. It is bounded by maximum flooding surface at the base and composite surface at the top.
Regressive systems tract forms in the marine part of the basin during the base level fall. Subaerial unconformities form in the landward side of the basin at the same time.
Parasequences and stacking patterns
A parasequence is a relatively conformable, genetically related succession of beds and bedsets bounded by marine flooding surfaces and their correlative surfaces. The flooding surfaces bounding parasequences are not of the same scale as the regional transgressive surface that is associated with a sequence boundary.
The parasequences are separated into stacking patterns:
Aggradational
Progradational
Retrogradational
Each stacking pattern will give different information on the behaviour of accommodation space, a major control of which is relative level. So a rapidly progradational pattern will be indicative of falling sea level, rapidly retrogradational is evidence for rapidly transgressing sea level and aggradational will be indicative of gently rising sea level.
Sea level through geologic time
Sea level changes over geologic time. The graph on the right illustrates two recent interpretations of sea level changes during the Phanerozoic. The modern age is depicted on the left side, labeled N for Neogene.
The blue spikes near date zero represent the sea level changes associated with the most recent glacial period, which reached its maximum extent about 20,000 years Before Present (BP). During this glaciation event, the world's sea level was about 320 feet (98 meters) lower than today, due to the large amount of sea water that had evaporated and been deposited as snow and ice in Northern Hemisphere glaciers. When the world's sea level was at this "low stand", former sea bed sediments were subjected to subaerial weathering (erosion by rain, frost, rivers, etc.) and a new shoreline was established at the new level, sometimes miles basinward of the former shoreline if the sea floor was shallowly inclined.
Today, sea level is at a relative "high stand" within the Quaternary glacial cycles because of rapid end-Pleistocene and early-Holocene deglaciation. The ancient shoreline of the last glacial period is now under approximately of water. Although there is debate among earth scientists whether we are currently experiencing a "high stand" it is generally accepted that the eustatic sea level is rising.
In the distant past, sea level has been significantly higher than today. During the Cretaceous (labeled K on the graph), sea level was so high that a seaway extended across the center of North America from Texas to the Arctic Ocean.
These alternating high and low sea level stands repeat at several time scales. The smallest of these cycles is approximately 20,000 years, and corresponds to the rate of precession of the Earth's rotational axis (see Milankovitch cycles) and are commonly referred to as '5th order' cycles. The next larger cycle ('4th order') is about 40,000 years and approximately matches the rate at which the Earth's inclination to the Sun varies (again explained by Milankovitch). The next larger cycle ('3rd order') is about 110,000 years and corresponds to the rate at which the Earth's orbit oscillates from elliptical to circular. Lower order cycles are recognized, which seem to result from plate tectonic events like the opening of new ocean basins by splitting continental masses.
Hundreds of similar glacial cycles have occurred throughout the Earth's history. The earth scientists who study the positions of coastal sediment deposits through time ("sequence stratigraphers") have noted dozens of similar basinward shifts of shorelines associated with a later recovery. The largest of these sedimentary cycles can in some cases be correlated around the world with great confidence.
The three controls on stratigraphic architecture and sedimentary cycle development are:
Eustatic sea level changes
Subsidence rate of the basin
Sediment supply.
Eustatic sea level is the sea level with reference to a fixed point, the centre of the Earth. Relative sea level is measured with reference to the base level, above which erosion can occur and below which deposition can occur. Both eustatic sea level changes and subsidence rates tend to be longer cycles. Sediment supply is largely thought to be controlled by local climatic conditions and can vary rapidly. These variations in local sediment supply affect the local and relative sea level which causes local sedimentary cycles.
Smaller and localised sedimentary cycles are not related to worldwide (eustatic) sea level changes but more to the supply of sediment to the adjacent basins where these sediments are being supplied. For example, when the basinward (oceanward) shift with progradation of shorelines was occurring in the Book Cliffs area of Utah the shorelines were receding or transgressing northwards in Wyoming. These sedimentary cycles are representative of the amount of supply of sediment to the basin. In a transgression, less sediment is being supplied than the rate of increase in the depth of water, and thus the shoreline migrates landward. In a regression, if the water depth is decreasing, the shoreline migrates seaward (basinward) and the previous shoreline is eroded. A regression of the shoreline also occurs if more sediment is being supplied than the shoreline can erode, causing the shoreline to migrate seaward. The latter is called progradation. The cycle of strata deposited during repeated transgressions and regressions creates a depositional sequence.
Economic significance
Sequence stratigraphy is and essential tool in the application of geology to the exploration for oil and gas, as a part of the field of Petroleum geology. Much of the development of this scientific discipline has occurred within or been funded by energy corporations and their geological research labs.
Sequence boundaries have economic significance because these changes in sea level cause large lateral shifts in the depositional patterns of seafloor sediments. These lateral shifts in deposition create alternating layers of good reservoir quality rock (porous and permeable sands) and poorer-quality mudstones (capable of providing a reservoir "seal" to prevent the leakage of any accumulated hydrocarbons that may have migrated into the sandstones). Hydrocarbon prospectors look for places in the world where porous and permeable sands are overlain by low permeability rocks, and where conditions are right for hydrocarbons to be generated and migrate into these "traps".
| Physical sciences | Stratigraphy | Earth science |
24556191 | https://en.wikipedia.org/wiki/European%20rabbit | European rabbit | The European rabbit (Oryctolagus cuniculus) or coney is a species of rabbit native to the Iberian Peninsula (Spain, Portugal and Andorra) and southwestern France. It is the only extant species in the genus Oryctolagus. The European rabbit has faced a population decline in its native range due to myxomatosis, rabbit hemorrhagic disease, overhunting and habitat loss. Outside of its native range, it is known as an invasive species, as it has been introduced to countries on all continents with the exception of Antarctica, often with devastating effects on local biodiversity due to a lack of predators.
The average adult European rabbit is in length, and can weigh , though size and weight vary with habitat and diet. Its distinctive ears can measure up to from the occiput. Due to the European rabbit's history of domestication, selective breeding, and introduction to non-native habitats, feral European rabbits across the world display a wide variety of morphologies.
The European rabbit is well known for digging networks of burrows, called warrens, where it spends most of its time when not feeding. It is a gregarious species, and lives in social groups centered around territorial females. European rabbits in an established social group will rarely stray far from their warren, with female rabbits leaving the warren mainly to establish nests where they will raise their young. Unlike hares (Lepus spp.), rabbits are altricial and are born blind, requiring maternal care until they leave the nest after 18 days.
Much of the modern research into wild rabbit behaviour was carried out in the 1960s by two research centres. One was the naturalist Ronald Lockley, who maintained a number of large enclosures for wild rabbit colonies, with observation facilities at Orielton, in Pembrokeshire, Wales. Apart from publishing a number of scientific papers, he popularised his findings in a book The Private Life of the Rabbit, which is credited by Richard Adams as having played a key role in his gaining "a knowledge of rabbits and their ways" that informed his novel Watership Down. The other group was the Commonwealth Scientific and Industrial Research Organisation (CSIRO) in Australia, where numerous studies of the social behaviour of wild rabbits were performed. Since the onset of myxomatosis, and the decline of the significance of the rabbit as an agricultural pest, few large-scale studies have been performed and many aspects of rabbit behaviour are still poorly understood.
Naming and etymology
Because of its non-British origin, the species does not have native names in English or Celtic, with the usual terms "cony" and "rabbit" being foreign loanwords. "Rabbit" is also pronounced as rabbidge, rabbert (North Devon) and rappit (Cheshire and Lancashire). More archaic spellings include rabbette (15th–16th centuries), rabet (15th–17th centuries), rabbet (16th-18th centuries), rabatte (16th century), rabytt (17th century) and rabit (18th century). The root word is the Walloon , which was once commonly used in Liège. itself is derived from the Middle Dutch , with the addition of the suffix -ett.
The term "cony" or "coney" antedates "rabbit", and first occurred during the 13th century to refer to the animal's pelt. Later, "cony" referred to the adult animal, while "rabbit" referred to the young. The root of "cony" is the old French or , of which the Norman was , plural or . Connil comes from the Latin . Its forerunner is the Greek (). The origin of itself is unclear: Ælian, who lived during the third century, linked the word to Celtiberian and later authors relate it to its Basque name unchi; Varo and Pliny connected it to cuneus, which refers to a wedge, thus making reference to the animal's digging ability. Later study of the etymology of has attested to its origin as a diminutive or adjectival form of the root word for "dog" () in Celtiberian.
The species' dwelling place is termed a warren or cony-garth. "Warren" comes from the Old English , itself derived from the Old French , , or . The root word is the Low Latin , which originally signified a preserve in general, only to be later used to refer specifically to an enclosure set apart for rabbits and hares. "Cony-garth" derives from the Middle English , which may be a compound of connynge+erthe (cony+earth). The term stems from the Old French or , and later . The root word is the Low Latin , the feminine form of the adjective , which pertains to the rabbit (as in the specific name, ). The generic name, Oryctolagus, derives from (, "burrowing") and (, "hare").
Taxonomy
Originally assigned to the genus Lepus, the European rabbit was consigned to its own genus in 1874 on account of its altricial young, its burrowing habits, and numerous skeletal characters. It is superficially similar to the North American cottontails (Sylvilagus) in that they are born blind and naked, have white flesh, and little sexual dimorphism. However, they differ in skull characteristics, and cottontails do not construct their own burrows as the European rabbit does. Molecular studies confirm that the resemblance between the two is due to convergent evolution, and that the European rabbit's closest relatives are the hispid hare (Caprolagus hispidus), the riverine rabbit (Bunolagus monticularis), and the Amami rabbit (Pentalagus furnessi).
The following cladogram is from Matthee et al., 2004, based on nuclear and mitochondrial gene analysis.
Subspecies
In 2005, six subspecies were recognised in Mammal Species of the World:
Genetic studies undertaken in 2008, however, indicate only two extant subspecies, O. c. algirus and O. c. cuniculus, native to the Iberian Peninsula, where most of the European rabbit's evolutionary history is centered; as of 2023, only these two subspecies are recognized. O. c. algirus and O. c. cuniculus occupy the south-west and north-east regions of the peninsula, respectively, naturally coming in contact in a region that spans the north-west to south-east, and likely diverged during the Quaternary glaciation 2 million years ago. Subspecies other than O. c. algirus and O. c. cuniculus have been recommended for abandonment, as they have very little evolutionary history and genetic diversity, and are likely not indigenous to the regions they occupy. Populations considered native to North Africa, such as those considered part of the subspecies O. c. habetensis, were likely introduced by Phonecians navigating the Mediterranean Sea; they are considered to be O. c. cuniculus, as are most other populations in regions the European rabbit was introduced to.
Fossil record
The oldest known fossils of the currently living European rabbit species, Oryctolagus cuniculus, appeared in the Middle Pleistocene age in southern Spain. The first identifiable fossils from Oryctolagus appeared during the Miocene epoch, and species such as O. laynensis, a presumed ancestor of O. cuniculus, and O. lacosti were recorded from 3.5 Mya up until the appearance of O. cuniculus 0.6 Mya. However, the European rabbit was the only member of its genus to survive to the Late Pleistocene, whereupon it spread from the Iberian Peninsula to the Mediterranean and northern Europe. Palaeoichnological evidence exists of European rabbits burrowing in and disturbing what are likely Neanderthal burial sites.
Description
The European rabbit is smaller than the European hare and mountain hare, and lacks black ear tips, as well as having proportionately shorter legs. An adult European rabbit can measure in length, and weigh . The hind foot measures in length, while the ears are long from the occiput.
Size and weight vary according to food and habitat quality, with rabbits living on light soil with nothing but grass to feed on being noticeably smaller than specimens living on highly cultivated farmlands with plenty of roots and clover. Pure European rabbits weighing and upwards are uncommon, but are occasionally reported. One large specimen, caught in February 1890 in Lichfield, was weighed at . Unlike the brown hare, the male European rabbit is more heavily built than the female. The penis is short, and lacks a baculum and true glans; the testicles, which are located in scrotal sacs to each side of the penis, can be retracted into the abdomen when food is scarce or when sexually inactive. Rudimentary nipples are also present in male rabbits.
The fur of the European rabbit is made up of soft down hair covered by stiff guard hairs, and is generally greyish-brown, though this is subject to much variation. The guard hairs are banded brown and black, or grey, while the nape of the neck and scrotum are reddish. The chest patch is brown, while the rest of the underparts are white or grey. A white star shape is often present on kits' foreheads, but rarely occurs in adults. The whiskers are long and black, and the feet are fully furred and buff-coloured. The tail has a white underside, which becomes prominent when escaping danger. This may act as a signal for other rabbits to run.
Moulting occurs once a year, beginning in March on the face and spreading over the back. The underfur is completely replaced by October–November. The European rabbit exhibits great variation in colour, from light sandy, to dark grey and completely black. Such variation depends largely on the amount of guard hairs relative to regular pelage. Melanists are not uncommon in mainland Europe, though albinoes are rare.
The skeleton and musculature of the European rabbit, like other leporids, is suited to survival by rapid escape from predators. The hind limbs are an exaggerated feature, being much longer and capable of producing more force than the forelimbs; their growth and use is correlated to that of the rest of the rabbit's body, as action pressure from the muscles creates force that is then distributed through the skeletal structures. Underuse of the rabbit's muscles leads to osteoporosis via bone rarefaction. The skull of the European rabbit displays a significant facial tilt of roughly 45° forward relative to the basicranium at rest, which supports their locomotion being mainly jumping or hopping (saltorial) rather than running (cursorial).
Life history and behaviour
Social and territorial behaviours
The European rabbit lives in warrens that contain 2–10 other individuals living in smaller groups to ensure greater breeding success. Territoriality and aggression contribute greatly to the rabbits' maturation process, and help ensure survival of the population. Females tend to be more territorial than males, although the areas most frequented by females are not defended. Territories are marked with dung hills. The size of the species' home range varies according to habitat, food, shelter, cover from predators, and breeding sites, though it is generally small, encompassing about . Except during times of low rabbit density and abundance of high-quality food, male ranges tend to be larger than those held by females. The European rabbit rarely strays far from its burrow; when feeding on cultivated fields, it typically only moves away from its burrow, and rarely . It may, however, move as far as after an abrupt change in environment, such as a harvest. This behaviour may be an antipredator adaptation, as rabbits in areas where predators are under rigorous control may move three times further from their burrows than those in areas without predator management.
The European rabbit is a gregarious animal, which lives in stable social groups centred around females and sharing access to one or more burrow systems. Social structures tend to be looser in areas where burrow construction is relatively easy. Dominance hierarchies exist in parallel for both bucks and does. Among bucks, status is determined through access to does, with dominant bucks siring the majority of the colony's offspring. The dominant does have priority access to the best nesting sites, with competition over such sites often leading to serious injury or death. Subordinate does, particularly in large colonies, typically resort to using single-entrance breeding spots far from the main warren, or may abandon the warren entirely, thus making themselves vulnerable to fox or badger predation.
Reproduction and development
In the European rabbit's mating system, dominant bucks exhibit polygyny, whereas lower-status individuals (both bucks and does) often form monogamous breeding relationships . Rabbits signal their readiness to copulate by marking other animals and inanimate objects with an odoriferous substance secreted though a chin gland, in a process known as "chinning". Though male European rabbits may sometimes be amicable with one another, fierce fights can erupt among bucks during the breeding season, which typically starts in autumn and continues through to spring. Occasionally, the mating season will extend into the summer. Introduced populations in the Southern Hemisphere experience breeding seasons during the other half of the calendar year. A succession of 4 to 5 litters (usually three to seven kittens each, on average five) are produced annually, but in overpopulated areas, pregnant does may lose all their embryos through intrauterine resorption. Shortly before giving birth, the doe constructs a separate burrow known as a "stop" or "stab", generally in an open field away from the main warren. These breeding burrows are typically a few feet long and are lined with grass and moss, as well as fur plucked from the doe's belly. The breeding burrow protects the kits from adult bucks and predators.
The gestation period of the European rabbit is 30 days, with the sex ratio of male to female kits tending to be 1:1. Greater maternal investment over male offspring may result in higher birth weights for bucks. Kits born to the dominant buck and doe—which enjoy better nesting and feeding grounds—tend to grow larger and stronger and become more dominant than those born to subordinate rabbits. Not uncommonly, European rabbits mate again immediately after giving birth, with some specimens having been observed to nurse previous young whilst pregnant. Female European rabbits may become pregnant at three months of age, but do not reach their full reproductive ability until they are two years old, after which they remain able to reproduce for 4 more years.
Female European rabbits nurse their kits once a night, for only a few minutes. After suckling is complete, the doe seals the entrance to the stop with soil and vegetation. In its native Iberian and southern French range, European rabbit young have a growth rate of per day, though such kittens in non-native ranges may grow per day. Weight at birth is and increases to by 21–25 days, during the weaning period. European rabbit kits are born blind, deaf, and nearly naked. The ears do not gain the power of motion until 10 days of age, and can be erected after 13. The eyes open 11 days after birth. At 18 days, the kittens begin to leave the burrow. Sexual maturity in bucks is attained at 4 months, while does can begin to breed at 3–5 months.
Dewlaps
A dewlap is a longitudinal flap of skin or similar flesh that hangs beneath the lower jaw or neck. It is a secondary sex characteristic in rabbits, caused by the presence of female sex hormones. They develop with puberty. A female rabbit who has been neutered before reaching sexual maturity will not develop a dewlap, and even if a doe is neutered after developing a dewlap, the dewlap will gradually disappear over several months. This also aligns with the results of injecting male rabbits with female sex hormones, specifically the ones from pregnant women's urine. The male rabbits developed dewlaps, which then gradually disappeared once administration had ceased. While it is unclear exactly what function a dewlap performs, pregnant female rabbits will pluck fur from their dewlaps shortly before giving birth to line a nest for their young.
Burrowing behaviour
The European rabbit's burrows occur mostly on slopes and banks, where drainage is more efficient. The burrow entrances are typically in diameter, and are easily recognisable by the bare earth at their mouths. Vegetation growth is prevented by the constant passing and repassing of the resident rabbits. Big burrows are complex excavations which may descend to depths of several feet. They are not constructed on any specified plan, and appear to be enlarged or improved as a result of the promiscuous activity of several generations. Digging is done by pulling the soil backwards with the fore feet and throwing it between the hind legs, which scatter the material with kicking motions. While most burrows are dug from the outside, some warrens feature holes dug from the inside, which act as emergency exits when escaping from predators below ground. These holes usually descend perpendicularly to , and their mouths lack the bare-earth characteristic of burrow entrances. While kits sleep in chambers lined with grass and fur, adults sleep on the bare earth, likely to escape dampness, with warmth being secured by huddling. Although both sexes dig, does do so more skillfully, and for longer periods.
Communication
The European rabbit is a relatively quiet animal, though it has at least two vocalisations. The best-known is a high treble scream or squeal. This distress call has been likened to the cry of a piglet. This sound is uttered when in extreme distress, such as being caught by a predator or trap. During the spring, bucks express contentment by emitting grunting sounds when approaching other rabbits. These grunts are similar to shrill hiccups, and are emitted with the mouth closed. Aggression is expressed with a low growl.
Ecology
Habitat
The European rabbit's ideal habitat consists of short grasslands with secure refuge (such as burrows, boulders, hedgerows, scrub, and woodland) near feeding areas. It may dwell up to treeline, as long as the land is well-drained and shelter is available. The size and distribution of its burrow systems depend on the type of soil present. In areas with loose soil, it selects sites with supporting structures, such as tree roots or shrubs to prevent burrow collapse. Warrens tend to be larger and have more interconnected tunnels in areas with chalk than those in sand. In large coniferous plantations, the species only occurs on peripheral areas and along fire breaks and rides. The European rabbit's grazing habits tend to promote their ideal open grassland habitat via the dispersion of seeds and trimming of vegetation.
Diet
The European rabbit eats a wide variety of herbage, especially grasses, favouring the young, succulent leaves and shoots of the most nutritious species, particularly fescues. In mixed cultivated areas, winter wheat is preferred over maize and dicotyledons. During the summer, the European rabbit feeds on the shortest, and therefore less nutritious grass swards, thus indicating that grazing grounds are selected through antipredator considerations rather than maximising food intake. In times of scarcity, the rabbit increases its food intake, selecting the parts of the plant with the highest nitrogen content. Hungry rabbits in winter may resort to eating tree bark. Blackberries are also eaten, and captive-bred European rabbits have been fed on fodder consisting of furze and acorns, which can lead to considerable weight gain. The European rabbit is a less fussy eater than the brown hare. When eating root vegetables, the rabbit eats them whole, while the hare tends to leave the peel. Depending on the body's fat and protein reserves, the species can survive without food in winter for about 2–8 days. Although herbivorous, cases are known of rabbits eating snails.
Like other leporids, the European rabbit produces soft, mucus-covered faecal pellets, which are ingested directly from the anus. The soft pellets are produced posterior to the colon in the hind gut soon after the excretion of hard pellets and the stomach begins to fill with newly grazed food. The soft pellets are filled with protein-rich bacteria, and pass down to the rectum in glossy clusters. The rabbit swallows them whole, without perforating the enveloping membrane.
Predators
The European rabbit is prey to many different predatory species. Foxes, dingoes, wolves, lynxes, wolverines, and dogs kill both adult and young rabbits by stalking and surprising them in the open, but relatively few rabbits are caught this way, as they can quickly rush back to cover with a burst of speed. Further, evidence from a study in Spain suggests they may avoid areas where the recent scat of predators which have eaten rabbit is detected. Both foxes and badgers dig out kittens from shallow burrows, with the latter predators being too slow to catch adult rabbits. Both wild and domestic cats can stalk and leap upon rabbits, particularly young specimens leaving their burrows for the first time. Wildcats take rabbits according to availability; in eastern Scotland, where rabbits are abundant, they can make up over 90% of the wildcats' diet. Most domestic cats are incapable of killing healthy, full-grown adults, but will take weak and diseased ones. Does can be fiercely protective of their kits, having been observed to chase away large cats and mustelids, including ferrets, stoats, and weasels. However, rabbits typically run from mustelids, and may fear them innately. Cases are known of rabbits becoming paralysed with fear and dying when pursued by stoats or weasels, even when rescued unharmed.
The European rabbit makes up 85% of the polecat's diet, and its availability is important to the success of breeding female mink. Brown rats can be a serious threat to kittens, as they will reside in rabbit burrows during the summer, and attack them in groups. Although many birds of prey are capable of killing rabbits, few are strong enough to carry them. Large species, such as golden and sea eagles, may carry rabbits back to their nests, while small eagles, buzzards, and harriers struggle to do so. Hawks and owls typically only carry off very small kits. Due to its decline in the Iberian Peninsula, the Iberian lynx (Lynx pardinus) and Spanish imperial eagle (Aquila adalberti), specialist predators of the European rabbit, have faced subsequent downturns in population.
Diseases, parasites and immunity
The European rabbit is the only species fatally attacked by myxomatosis. The most lethal strain has a five-day incubation period, after which the eyelids swell, with the inflammation quickly spreading to the base of the ears, the forehead, and nose. At the same time, the anal and genital areas also swell. During the last stages of the disease, the swellings discharge a fluid rich in viral material, with death usually following on the 11th–12th day of infection. The primary carrier of myxomatosis varies based on location; in North America and Australia, it is carried by multiple species of mosquitoes, while in Britain its primary carrier is the rabbit flea (Spilopsyllus cuniculi).
Rabbit haemorrhagic disease (RHD), also known as viral haemorrhagic disease or rabbit calicivirus disease in Australia, is specific to the European rabbit, and causes lesions of acute necrotising hepatitis, disseminated intravascular coagulation, and haemorrhaging, mainly in the lungs. Susceptible specimens may die within 30 hours of infection. Most rabbits in the UK are immune to RHD, due to exposure to a weaker strain.
European rabbit immunity has significantly diverged from other tetrapods in the manner in which it employs immunoglobulin light chains. In one case, McCartney-Francis (et al., 1984) discovered a unique additional disulfide bond between Cys 80 in Vκ and Cys 171 in Cκ. They suggest that this may serve to stabilise rabbit antibodies. The gene , responsible for the principal immunoglobulin light chain, shows high amino acid divergence between domesticated types and ferals derived from them. This divergence can be as high as 40%, and indicates high genetic diversity of populations surviving over evolutionary time scales.
Human relationships with rabbits
Recent research has shown that all European rabbits carry common genetic markers and descend from one of two maternal lines. These lines originated between 12,000 and 6.5 million years ago when glaciers isolated two herds, one on the Iberian Peninsula and the other in Southern France. Humans likely began hunting rabbits as a food source, but further research is needed to verify this. Little comprehensive evidence of the relationship of humans with European rabbits is documented until the medieval period.
Humans' relationship with the European rabbit was first recorded by the Phoenicians prior to 1000 BC, when they termed the Iberian Peninsula i-Shaphan-ím (literally, the land of the hyraxes). This phrase closely resembles related modern Hebrew: I (אי) meaning island and shafan (שפן) meaning hyrax, plural shfaním (שפנים). Phoenicians called the local rabbits 'hyraxes' because rabbits resemble hyraxes in some ways, and hyraxes are native to Phoenicia, unlike rabbits. Hyraxes, like rabbits, are not rodents. One theory states that the Romans converted the phrase i-Shaphan-ím, with influence from the Greek Spania, to its Latin form, Hispania, which evolved in all the Iberian languages - into Castilian España, Portuguese Espanha, Catalan Espanya (English "Spain"), and such other variations in modern languages. Different views have been voiced on the precise meaning of shafan, but the balance of opinion appears to indicate that the hyrax is indeed the intended meaning.
Like the Phoenicians, neither the later Greek nor Roman colonizers had a specific name for the rabbit, because the species wasn't native to Greece and Italy (though it is present there nowadays). They commonly called it "small hare" and "small digging hare", in contrast to the European hare, which is larger and does not make burrows. Catullus used the name cuniculus (a latinization of the Western Iberian word and the etymological origin of the Castilian name conejo, Portuguese coelho and Catalan conill, and the English name coney), and referenced its abundance in Celtiberia by calling this region cuniculosa, i.e. rabbit-ridden.
The European rabbit is the only rabbit species that has been domesticated and all 305 global rabbit breeds— from Netherland Dwarf to Flemish Giant— are descendants of the European rabbit. Rabbits are an example of an animal that can be treated as a food, a pet, or a pest by different members of the same culture. In some urban areas, infestations of feral European rabbits (descended from pets) have become a problem. Helsinki, for example, host to one of the northernmost populations of the species, had an estimated 2,500 European rabbits at the end of 2006, doubling to 5,000 by autumn 2007.
As an introduced species
The European rabbit has been introduced as an exotic species into several environments, often with harmful results to vegetation and local wildlife, making it an invasive species. The first known mention of the rabbit as an invasive species (and possibly the first documented instance of an invasive species ever) was made in regard to the introduction of the rabbit to the Balearic Islands after the Roman conquest of the first century BCE. According to both Strabo and Pliny the Elder, the multiplying rabbits caused famines by destroying crop yields and even collapsed trees and houses with their burrowing. The inhabitants petitioned Augustus for help, who sent troops to curb the rabbit population with the help of ferrets.
Other locations where the European rabbit was introduced include Great Britain; two of the Hawaiian Islands (Laysan Island and Lisianski Island); Oceania's Macquarie Island; the island Ōkunoshima in Japan; Washington's Smith Island and San Juan Island (around 1900 and later spreading to the other San Juan Islands); several islands off the coast of Southern Africa (including Robben Island); and Australia and New Zealand. The two accounts over the introduction of rabbits in Ukraine are conflicting. One holds that the species was brought there in the early 20th century by Austrian nobleman Graf Malokhovsky, who released them on his estate near the Khadzhibey Estuary, while another holds that rabbits were first brought to Kherson from Switzerland in 1894–1895 by landowner Pinkovsky.
In the British Isles
The European rabbit is widespread in Great Britain, Ireland, and most other islands, except for Isles of Scilly, Rùm, Tiree, and some small Scottish islands, such as Gunna, Sanday, and most of the Treshnish Isles. It was likely first brought to Britain by the Normans after the 1066 conquest of England, as no pre-Norman British allusions to the animal have been found. The rabbit was nonetheless scarce or absent throughout most of England a short time afterwards, as warrens are not mentioned in the Domesday Book or any other 11th–century documents. Rabbits became well known, but not necessarily accepted members of British fauna between the 12th and 13th centuries. The first real evidence of their presence is a number of bones from the midden of Rayleigh Castle, which was occupied from the 11th–13th centuries. The first references to rabbits in Ireland occur roughly at the same time as English ones, thus indicating another Norman introduction. They had become plentiful, probably at a local level, by the 13th century, as indicated by an inquisition of Lundy Island made in 1274 describing how 2,000 rabbits were caught annually. Subsequent allusions in official documents became more frequent, with the species later becoming an important food item at feasts.
Truly wild populations increased slowly, primarily in the coastal areas and lowland heaths of Breckland and Norfolk. There were notable population increases after 1750, when changes in agricultural practices created favourable habitats, and increasing interest in game management resulted in intensive predator control campaigns. Although now common in the Scottish lowlands, the species was little known in Scotland before the 19th century. Until then, it was confined to portions of the Edinburgh district at least as far back as the 16th century, certain islands and the coastal sand dunes of the Scottish mainland. Although unknown in Caithness in 1743, the species became well established there by 1793.
Myxomatosis entered Britain from France in 1953, and reached Ireland by 1954, prompting the RSPB to set up "mercy squads" meant to euthanise myxomatous rabbits. Major myxomatosis outbreaks still occur in Britain, peaking twice annually: in spring and especially in late summer or autumn, though immunity has reduced the mortality rate from 99% to 5–33%.
Between 1996 and 2018, rabbit numbers fell by 88% in the east Midlands of England, 83% in Scotland, and 43% across the whole of the UK. Numbers are still falling (in 2021). Pip Mountjoy, Shifting Sands project manager at Natural England, said: "They (rabbits) are actually an endangered species in their native region on the Iberian peninsula. It is surprising for people that rabbits are important in some ecosystems. We think of them as a pest but in Britain they are a keystone species – they act as landscape managers and a lot of other species rely on them." The Shifting Sands project aims to encourage landowners to create safe habitats for rabbits, consisting of piles of branches placed near existing rabbit warrens. Species that depend on rabbits' grazing habits include purple milk vetch, rare spring sedge, spring speedwell, prostrate perennial knawel, caterpillars of the lunar yellow underwing moth, stone curlew, and the large blue butterfly.
In Australia
Though rabbits were first introduced to Australia in 1788 with the arrival of the First Fleet, the most significant population explosion occurred later on in the 19th century. Twenty-four specimens of the European rabbit were introduced to Australia in 1859 by estate owner Thomas Austin in Victoria. Their descendants multiplied and spread throughout the country and caused severe agricultural damage and widespread ecological changes that contributed to the decline of native Austrlian species such as the greater bilby (Macrotis lagotis) and the southern pig-footed bandicoot (Chaeropus ecaudatus).
Between 1901 and 1907, Australia built an immense "rabbit-proof fence" to halt the westward expansion of the infestation. The European rabbit, however, can not only jump very high, but also burrow underground, and this fence failed to protect from rabbit infestation; despite this, further fencing projects were undertaken that also failed to control the spread of rabbits in Australia.
During the 1950s, the intentional introduction of a virus that causes myxomatosis provided some relief in Australia, but not in New Zealand, where the insect vectors necessary for the spread of the disease were not present. Myxomatosis can also infect pet rabbits (the same species). Today's remaining feral rabbits in Australia are largely immune to myxomatosis. A strain of a second deadly rabbit virus, rabbit hemorrhagic disease (RHD), was imported to Australia in 1991 as a biological control agent, and was released accidentally in 1995, killing millions of rabbits. The virus has been developed further to address changes in environment and population. RHD was also introduced—illegally—in New Zealand with less success due to improper timing.
In Chile
The exact date on which the European rabbit was introduced into Chile is unknown, though the first references to it occur during the mid-18th century. By the 19th century, several authors referred to the presence of both rabbits and rabbit hutches in central Chile. The importation and breeding of rabbits was encouraged by the state, as rabbits were seen as cheap sources of food for peasants. Whether or not their escape into the wild was intentional is unknown, but warnings over the dangers of feral rabbits were raised during the early 20th century, and the species had propagated dramatically by the late 1920s in central Chile, Tierra del Fuego, and the Juan Fernández Islands. In the 1930s, the state sought to tackle the rabbit problem by banning fox hunting, though it was later discovered that indigenous South American foxes rarely preyed on rabbits, preferring native species. In modern times, the European rabbit problem has not been resolved definitively, though a deliberate outbreak of myxomatosis in Tierra del Fuego successfully reduced local rabbit populations. The species remains a problem in central Chile and on Juan Fernández, despite international financing.
Domestication
The European rabbit is the only rabbit to be widely domesticated, for meat, fur, wool, or as a pet. It was first widely kept in ancient Rome from the first century BC, where Pliny the Elder described the use of rabbit hutches, along with enclosures called leporaria. The European rabbit has been refined into a wide variety of breeds during and since the emergence of animal fancy in the 19th century.
Selective breeding has been used since ancient times in efforts to raise rabbits with different characteristics, and while domestic rabbits are typically larger than wild rabbits, the various breeds of domestic rabbit exist in a range of sizes from "dwarf" to "giant". They have as much color variation among themselves as other livestock and pet animals. Their fur is prized for its softness; the large Angora rabbit breeds are raised for their long, soft fur, which is often spun into yarn. Other breeds are raised for the fur industry, particularly the Rex, which has a smooth, velvet-like coat.
Meat and fur
In the United Kingdom, rabbit was a popular food source for the poorer classes. Among wild rabbits, those native to Spain were reputed to have the highest meat quality, followed by those in the Ardennes. As rabbits hold very little fat, they were hardly ever roasted, being instead boiled, fried, or stewed.
The pelt of the rabbit is heavier and more durable than the hare's. Marshall calculated that the value of the skin in proportion to the carcass was greater than that of the sheep and ox. Its fur is primarily used for felting or hats. It is also dyed or clipped, and sold as imitations of more valuable furbearers, such as fur seal. Although cheap and easily acquired, rabbit fur has little durability.
Conservation status
Though the European rabbit thrives in many of the locations where it was introduced, in its native Iberia, populations are dwindling. In 2005, the Portuguese Institute for Nature Conservation and Forests classified O. cuniculus in Portugal as "near threatened", while in 2006, Spanish authorities (SECEM) reclassified it in Spain as "vulnerable". In 2018, the International Union for Conservation of Nature reclassified O. cuniculus in Spain, Portugal, and France as "endangered", due to the extent of recent declines. The IUCN assessment of the species considers only those populations within its natural distribution, and as such it is considered endangered by that group.
| Biology and health sciences | Lagomorphs | Animals |
40017873 | https://en.wikipedia.org/wiki/Diabetes | Diabetes | Diabetes, also known as diabetes mellitus, is a group of common endocrine diseases characterized by sustained high blood sugar levels. Diabetes is due to either the pancreas not producing enough insulin, or the cells of the body becoming unresponsive to the hormone's effects. Classic symptoms include polydipsia (excessive thirst), polyuria (excessive urination), weight loss, and blurred vision. If left untreated, the disease can lead to various health complications, including disorders of the cardiovascular system, eye, kidney, and nerves. Diabetes accounts for approximately 4.2 million deaths every year, with an estimated 1.5 million caused by either untreated or poorly treated diabetes.
The major types of diabetes are type 1 and type 2. The most common treatment for type 1 is insulin replacement therapy (insulin injections), while anti-diabetic medications (such as metformin and semaglutide) and lifestyle modifications can be used to manage type 2. Gestational diabetes, a form that arises during pregnancy in some women, normally resolves shortly after delivery.
The number of people diagnosed as living with diabetes has increased sharply in recent decades, from 200 million in 1990 to 830 million by 2022. It affects one in seven of the adult population, with type 2 diabetes accounting for more than 95% of cases. These numbers have already risen beyond earlier projections of 783 million adults by 2045. The prevalence of the disease continues to increase, most dramatically in low- and middle-income nations. Rates are similar in women and men, with diabetes being the seventh leading cause of death globally. The global expenditure on diabetes-related healthcare is an estimated US$760 billion a year.
Signs and symptoms
Common symptoms of diabetes include increased thirst, frequent urination, extreme hunger, and unintended weight loss. Several other non-specific signs and symptoms may also occur, including fatigue, blurred vision, sweet smelling urine/semen and genital itchiness due to Candida infection. About half of affected individuals may also be asymptomatic. Type 1 presents abruptly following a pre-clinical phase, while type 2 has a more insidious onset; patients may remain asymptomatic for many years.
Diabetic ketoacidosis is a medical emergency that occurs most commonly in type 1, but may also occur in type 2 if it has been longstanding or if the individual has significant β-cell dysfunction. Excessive production of ketone bodies leads to signs and symptoms including nausea, vomiting, abdominal pain, the smell of acetone in the breath, deep breathing known as Kussmaul breathing, and in severe cases decreased level of consciousness. Hyperosmolar hyperglycemic state is another emergency characterized by dehydration secondary to severe hyperglycemia, with resultant hypernatremia leading to an altered mental state and possibly coma.
Hypoglycemia is a recognized complication of insulin treatment used in diabetes. An acute presentation can include mild symptoms such as sweating, trembling, and palpitations, to more serious effects including impaired cognition, confusion, seizures, coma, and rarely death. Recurrent hypoglycemic episodes may lower the glycemic threshold at which symptoms occur, meaning mild symptoms may not appear before cognitive deterioration begins to occur.
Long-term complications
The major long-term complications of diabetes relate to damage to blood vessels at both macrovascular and microvascular levels. Diabetes doubles the risk of cardiovascular disease, and about 75% of deaths in people with diabetes are due to coronary artery disease. Other macrovascular morbidities include stroke and peripheral artery disease.
Microvascular disease affects the eyes, kidneys, and nerves. Damage to the retina, known as diabetic retinopathy, is the most common cause of blindness in people of working age. The eyes can also be affected in other ways, including development of cataract and glaucoma. It is recommended that people with diabetes visit an optometrist or ophthalmologist once a year.
Diabetic nephropathy is a major cause of chronic kidney disease, accounting for over 50% of patients on dialysis in the United States. Diabetic neuropathy, damage to nerves, manifests in various ways, including sensory loss, neuropathic pain, and autonomic dysfunction (such as postural hypotension, diarrhoea, and erectile dysfunction). Loss of pain sensation predisposes to trauma that can lead to diabetic foot problems (such as ulceration), the most common cause of non-traumatic lower-limb amputation.
Hearing loss is another long-term complication associated with diabetes.
Based on extensive data and numerous cases of gallstone disease, it appears that a causal link might exist between type 2 diabetes and gallstones. People with diabetes are at a higher risk of developing gallstones compared to those without diabetes.
There is a link between cognitive deficit and diabetes; studies have shown that diabetic individuals are at a greater risk of cognitive decline, and have a greater rate of decline compared to those without the disease. The condition also predisposes to falls in the elderly, especially those treated with insulin.
Types
Diabetes is classified by the World Health Organization into six categories: type 1 diabetes, type 2 diabetes, hybrid forms of diabetes (including slowly evolving, immune-mediated diabetes of adults and ketosis-prone type 2 diabetes), hyperglycemia first detected during pregnancy, "other specific types", and "unclassified diabetes". Diabetes is a more variable disease than once thought, and individuals may have a combination of forms.
Type 1
Type 1 accounts for 5 to 10% of diabetes cases and is the most common type diagnosed in patients under 20 years; however, the older term "juvenile-onset diabetes" is no longer used as onset in adulthood is not unusual. The disease is characterized by loss of the insulin-producing beta cells of the pancreatic islets, leading to severe insulin deficiency, and can be further classified as immune-mediated or idiopathic (without known cause). The majority of cases are immune-mediated, in which a T cell-mediated autoimmune attack causes loss of beta cells and thus insulin deficiency. Patients often have irregular and unpredictable blood sugar levels due to very low insulin and an impaired counter-response to hypoglycemia.
Type 1 diabetes is partly inherited, with multiple genes, including certain HLA genotypes, known to influence the risk of diabetes. In genetically susceptible people, the onset of diabetes can be triggered by one or more environmental factors, such as a viral infection or diet. Several viruses have been implicated, but to date there is no stringent evidence to support this hypothesis in humans.
Type 1 diabetes can occur at any age, and a significant proportion is diagnosed during adulthood. Latent autoimmune diabetes of adults (LADA) is the diagnostic term applied when type 1 diabetes develops in adults; it has a slower onset than the same condition in children. Given this difference, some use the unofficial term "type 1.5 diabetes" for this condition. Adults with LADA are frequently initially misdiagnosed as having type 2 diabetes, based on age rather than a cause. LADA leaves adults with higher levels of insulin production than type 1 diabetes, but not enough insulin production for healthy blood sugar levels.
Type 2
Type 2 diabetes is characterized by insulin resistance, which may be combined with relatively reduced insulin secretion. The defective responsiveness of body tissues to insulin is believed to involve the insulin receptor. However, the specific defects are not known. Diabetes mellitus cases due to a known defect are classified separately. Type 2 diabetes is the most common type of diabetes mellitus accounting for 95% of diabetes. Many people with type 2 diabetes have evidence of prediabetes (impaired fasting glucose and/or impaired glucose tolerance) before meeting the criteria for type 2 diabetes. The progression of prediabetes to overt type 2 diabetes can be slowed or reversed by lifestyle changes or medications that improve insulin sensitivity or reduce the liver's glucose production.
Type 2 diabetes is primarily due to lifestyle factors and genetics. A number of lifestyle factors are known to be important to the development of type 2 diabetes, including obesity (defined by a body mass index of greater than 30), lack of physical activity, poor diet such as Western Pattern Diet, stress, and urbanization. Excess body fat is associated with 30% of cases in people of Chinese and Japanese descent, 60–80% of cases in those of European and African descent, and 100% of Pima Indians and Pacific Islanders. Even those who are not obese may have a high waist–hip ratio.
Dietary factors such as sugar-sweetened drinks are associated with an increased risk. The type of fats in the diet is also important, with saturated fat and trans fats increasing the risk and polyunsaturated and monounsaturated fat decreasing the risk. Eating white rice excessively may increase the risk of diabetes, especially in Chinese and Japanese people. Lack of physical activity may increase the risk of diabetes in some people.
Adverse childhood experiences, including abuse, neglect, and household difficulties, increase the likelihood of type 2 diabetes later in life by 32%, with neglect having the strongest effect.
Antipsychotic medication side effects (specifically metabolic abnormalities, dyslipidemia and weight gain) are also potential risk factors.
Gestational diabetes
Gestational diabetes resembles type 2 diabetes in several respects, involving a combination of relatively inadequate insulin secretion and responsiveness. It occurs in about 2–10% of all pregnancies and may improve or disappear after delivery. It is recommended that all pregnant women get tested starting around 24–28 weeks gestation. It is most often diagnosed in the second or third trimester because of the increase in insulin-antagonist hormone levels that occurs at this time. However, after pregnancy approximately 5–10% of women with gestational diabetes are found to have another form of diabetes, most commonly type 2. Gestational diabetes is fully treatable, but requires careful medical supervision throughout the pregnancy. Management may include dietary changes, blood glucose monitoring, and in some cases, insulin may be required.
Though it may be transient, untreated gestational diabetes can damage the health of the fetus or mother. Risks to the baby include macrosomia (high birth weight), congenital heart and central nervous system abnormalities, and skeletal muscle malformations. Increased levels of insulin in a fetus's blood may inhibit fetal surfactant production and cause infant respiratory distress syndrome. A high blood bilirubin level may result from red blood cell destruction. In severe cases, perinatal death may occur, most commonly as a result of poor placental perfusion due to vascular impairment. Labor induction may be indicated with decreased placental function. A caesarean section may be performed if there is marked fetal distress or an increased risk of injury associated with macrosomia, such as shoulder dystocia.
Other types
Maturity onset diabetes of the young (MODY) is a rare autosomal dominant inherited form of diabetes, due to one of several single-gene mutations causing defects in insulin production. It is significantly less common than the three main types, constituting 1–2% of all cases. The name of this disease refers to early hypotheses as to its nature. Being due to a defective gene, this disease varies in age at presentation and in severity according to the specific gene defect; thus, there are at least 13 subtypes of MODY. People with MODY often can control it without using insulin.
Some cases of diabetes are caused by the body's tissue receptors not responding to insulin (even when insulin levels are normal, which is what separates it from type 2 diabetes); this form is very uncommon. Genetic mutations (autosomal or mitochondrial) can lead to defects in beta cell function. Abnormal insulin action may also have been genetically determined in some cases. Any disease that causes extensive damage to the pancreas may lead to diabetes (for example, chronic pancreatitis and cystic fibrosis). Diseases associated with excessive secretion of insulin-antagonistic hormones can cause diabetes (which is typically resolved once the hormone excess is removed). Many drugs impair insulin secretion and some toxins damage pancreatic beta cells, whereas others increase insulin resistance (especially glucocorticoids which can provoke "steroid diabetes"). The ICD-10 (1992) diagnostic entity, malnutrition-related diabetes mellitus (ICD-10 code E12), was deprecated by the World Health Organization (WHO) when the current taxonomy was introduced in 1999. Yet another form of diabetes that people may develop is double diabetes. This is when a type 1 diabetic becomes insulin resistant, the hallmark for type 2 diabetes or has a family history for type 2 diabetes. It was first discovered in 1990 or 1991.
The following is a list of disorders that may increase the risk of diabetes:
Genetic defects of β-cell function
Maturity onset diabetes of the young
Mitochondrial DNA mutations
Genetic defects in insulin processing or insulin action
Defects in proinsulin conversion
Insulin gene mutations
Insulin receptor mutations
Exocrine pancreatic defects (see Type 3c diabetes, i.e. pancreatogenic diabetes)
Chronic pancreatitis
Pancreatectomy
Pancreatic neoplasia
Cystic fibrosis
Hemochromatosis
Fibrocalculous pancreatopathy
Endocrinopathies
Growth hormone excess (acromegaly)
Cushing syndrome
Hyperthyroidism
Hypothyroidism
Pheochromocytoma
Glucagonoma
Infections
Cytomegalovirus infection
Coxsackievirus B
Drugs
Glucocorticoids
Thyroid hormone
β-adrenergic agonists
Statins
Pathophysiology
Insulin is the principal hormone that regulates the uptake of glucose from the blood into most cells of the body, especially liver, adipose tissue and muscle, except smooth muscle, in which insulin acts via the IGF-1. Therefore, deficiency of insulin or the insensitivity of its receptors play a central role in all forms of diabetes mellitus.
The body obtains glucose from three main sources: the intestinal absorption of food; the breakdown of glycogen (glycogenolysis), the storage form of glucose found in the liver; and gluconeogenesis, the generation of glucose from non-carbohydrate substrates in the body. Insulin plays a critical role in regulating glucose levels in the body. Insulin can inhibit the breakdown of glycogen or the process of gluconeogenesis, it can stimulate the transport of glucose into fat and muscle cells, and it can stimulate the storage of glucose in the form of glycogen.
Insulin is released into the blood by beta cells (β-cells), found in the islets of Langerhans in the pancreas, in response to rising levels of blood glucose, typically after eating. Insulin is used by about two-thirds of the body's cells to absorb glucose from the blood for use as fuel, for conversion to other needed molecules, or for storage. Lower glucose levels result in decreased insulin release from the beta cells and in the breakdown of glycogen to glucose. This process is mainly controlled by the hormone glucagon, which acts in the opposite manner to insulin.
If the amount of insulin available is insufficient, or if cells respond poorly to the effects of insulin (insulin resistance), or if the insulin itself is defective, then glucose is not absorbed properly by the body cells that require it, and is not stored appropriately in the liver and muscles. The net effect is persistently high levels of blood glucose, poor protein synthesis, and other metabolic derangements, such as metabolic acidosis in cases of complete insulin deficiency.
When there is too much glucose in the blood for a long time, the kidneys cannot absorb it all (reach a threshold of reabsorption) and the extra glucose gets passed out of the body through urine (glycosuria). This increases the osmotic pressure of the urine and inhibits reabsorption of water by the kidney, resulting in increased urine production (polyuria) and increased fluid loss. Lost blood volume is replaced osmotically from water in body cells and other body compartments, causing dehydration and increased thirst (polydipsia). In addition, intracellular glucose deficiency stimulates appetite leading to excessive food intake (polyphagia).
Diagnosis
Diabetes mellitus is diagnosed with a test for the glucose content in the blood, and is diagnosed by demonstrating any one of the following:
Fasting plasma glucose level ≥ 7.0 mmol/L (126 mg/dL). For this test, blood is taken after a period of fasting, i.e. in the morning before breakfast, after the patient had sufficient time to fast overnight or at least 8 hours before the test.
Plasma glucose ≥ 11.1 mmol/L (200 mg/dL) two hours after a 75 gram oral glucose load as in a glucose tolerance test (OGTT)
Symptoms of high blood sugar and plasma glucose ≥ 11.1 mmol/L (200 mg/dL) either while fasting or not fasting
Glycated hemoglobin (HbA1C) ≥ 48 mmol/mol (≥ 6.5 DCCT %).
A positive result, in the absence of unequivocal high blood sugar, should be confirmed by a repeat of any of the above methods on a different day. It is preferable to measure a fasting glucose level because of the ease of measurement and the considerable time commitment of formal glucose tolerance testing, which takes two hours to complete and offers no prognostic advantage over the fasting test. According to the current definition, two fasting glucose measurements at or above 7.0 mmol/L (126 mg/dL) is considered diagnostic for diabetes mellitus.
Per the WHO, people with fasting glucose levels from 6.1 to 6.9 mmol/L (110 to 125 mg/dL) are considered to have impaired fasting glucose. People with plasma glucose at or above 7.8 mmol/L (140 mg/dL), but not over 11.1 mmol/L (200 mg/dL), two hours after a 75 gram oral glucose load are considered to have impaired glucose tolerance. Of these two prediabetic states, the latter in particular is a major risk factor for progression to full-blown diabetes mellitus, as well as cardiovascular disease. The American Diabetes Association (ADA) since 2003 uses a slightly different range for impaired fasting glucose of 5.6 to 6.9 mmol/L (100 to 125 mg/dL).
Glycated hemoglobin is better than fasting glucose for determining risks of cardiovascular disease and death from any cause.
Prevention
There is no known preventive measure for type 1 diabetes. However, islet autoimmunity and multiple antibodies can be a strong predictor of the onset of type 1 diabetes. Type 2 diabetes—which accounts for 85–90% of all cases worldwide—can often be prevented or delayed by maintaining a normal body weight, engaging in physical activity, and eating a healthy diet. Higher levels of physical activity (more than 90 minutes per day) reduce the risk of diabetes by 28%. Dietary changes known to be effective in helping to prevent diabetes include maintaining a diet rich in whole grains and fiber, and choosing good fats, such as the polyunsaturated fats found in nuts, vegetable oils, and fish. Limiting sugary beverages and eating less red meat and other sources of saturated fat can also help prevent diabetes. Tobacco smoking is also associated with an increased risk of diabetes and its complications, so smoking cessation can be an important preventive measure as well.
The relationship between type 2 diabetes and the main modifiable risk factors (excess weight, unhealthy diet, physical inactivity and tobacco use) is similar in all regions of the world. There is growing evidence that the underlying determinants of diabetes are a reflection of the major forces driving social, economic and cultural change: globalization, urbanization, population aging, and the general health policy environment.
Comorbidity
Diabetes patients' comorbidities have a significant impact on medical expenses and related costs. It has been demonstrated that patients with diabetes are more likely to experience respiratory, urinary tract, and skin infections, develop atherosclerosis, hypertension, and chronic kidney disease, putting them at increased risk of infection and complications that require medical attention. Patients with diabetes mellitus are more likely to experience certain infections, such as COVID-19, with prevalence rates ranging from 5.3 to 35.5%. Maintaining adequate glycemic control is the primary goal of diabetes management since it is critical to managing diabetes and preventing or postponing such complications.
People with type 1 diabetes have higher rates of autoimmune disorders than the general population. An analysis of a type 1 diabetes registry found that 27% of the 25,000 participants had other autoimmune disorders. Between 2% and 16% of people with type 1 diabetes also have celiac disease.
Management
Diabetes management concentrates on keeping blood sugar levels close to normal, without causing low blood sugar. This can usually be accomplished with dietary changes, exercise, weight loss, and use of appropriate medications (insulin, oral medications).
Learning about the disease and actively participating in the treatment is important, since complications are far less common and less severe in people who have well-managed blood sugar levels. The goal of treatment is an A1C level below 7%. Attention is also paid to other health problems that may accelerate the negative effects of diabetes. These include smoking, high blood pressure, metabolic syndrome obesity, and lack of regular exercise. Specialized footwear is widely used to reduce the risk of diabetic foot ulcers by relieving the pressure on the foot. Foot examination for patients living with diabetes should be done annually which includes sensation testing, foot biomechanics, vascular integrity and foot structure.
Concerning those with severe mental illness, the efficacy of type 2 diabetes self-management interventions is still poorly explored, with insufficient scientific evidence to show whether these interventions have similar results to those observed in the general population.
Lifestyle
People with diabetes can benefit from education about the disease and treatment, dietary changes, and exercise, with the goal of keeping both short-term and long-term blood glucose levels within acceptable bounds. In addition, given the associated higher risks of cardiovascular disease, lifestyle modifications are recommended to control blood pressure.
Weight loss can prevent progression from prediabetes to diabetes type 2, decrease the risk of cardiovascular disease, or result in a partial remission in people with diabetes. No single dietary pattern is best for all people with diabetes. Healthy dietary patterns, such as the Mediterranean diet, low-carbohydrate diet, or DASH diet, are often recommended, although evidence does not support one over the others. According to the ADA, "reducing overall carbohydrate intake for individuals with diabetes has demonstrated the most evidence for improving glycemia", and for individuals with type 2 diabetes who cannot meet the glycemic targets or where reducing anti-glycemic medications is a priority, low or very-low carbohydrate diets are a viable approach. For overweight people with type 2 diabetes, any diet that achieves weight loss is effective.
A 2020 Cochrane systematic review compared several non-nutritive sweeteners to sugar, placebo and a nutritive low-calorie sweetener (tagatose), but the results were unclear for effects on HbA1c, body weight and adverse events. The studies included were mainly of very low-certainty and did not report on health-related quality of life, diabetes complications, all-cause mortality or socioeconomic effects.
Exercise has demonstrated to impact people’s lives for a better health outcome. However, fear of hypoglycemia can negatively impact exercise view on youth that have been diagnosed with diabetes. Managing insulin, carbohydrate intake, and physical activity becomes a task that drive youth away benefitting from enjoying exercises. With different studies, an understanding of what can be done and applied to the youth population diagnosed with Type 1 Diabetes has been conducted. A study’s aim was to focus on the impact of an exercise education on physical activity. During the length of a 12-month program, youth and their parents participated in 4 education sessions learning about the benefits, safe procedures, glucose control, and physical activity. With a survey conducted in the beginning, youth and parents demonstrated their fear of hypoglycemia. At the end of the program, most of the youth and parents showed confidence on how to manage and handle situations regarding hypoglycemia. In some instances, youth provided feedback that a continuation of the sessions would be beneficial. In two other studies, exercise was the aim to investigate on how it affects adolescents with T1D. In one of those studies, the impact was assessed in the changes of glucose in exercise by how many minutes per day, intensity, duration, and heart rate. Also, glucose was monitored to see changes during exercise, post exercise, and overnight. The other study investigated how types of exercises can affect glucose levels. The exercise types were continuous moderate exercise and interval-high-intensity exercise. Both types consisted of 2 sets of 10-minute work at different pedaling paces. The continuous pedaled at a 50% and had a 5-minute passive recovery. The high-intensity pedaled at 150% for 15 seconds and was intermixed with a 30-second passive recovery. So, when studies finished collecting data and were able to analyze it, the following were the results. For the studies comparing the different intensities, it was seen that insulin and carbohydrate intake did not have a significant difference before or after exercise. In regards of glucose content, there was a greater drop of blood glucose post exercise in the high intensity (-1.47mmol/L). During recovery, the continuous exercise showed a greater decrease in blood glucose. With all these, continuous exercise resulted in being more favorable for managing blood glucose levels. In the other study, it is mentioned that exercise also contributed to a notable impact on glucose levels. Post-exercise measurements, there was a low mean glucose level that occurred 12 to 16 hours after exercising. Although, with participants exercising for longer sessions (≥90 minutes), hypoglycemia rates were higher. With all these, participants showed well-managed glucose control by intaking proper carbohydrates amount without any insulin adjustments. Lastly, the study, that educated youth and parents about exercise important and management of hypoglycemia, showed many youths feeling confident to continue to exercise regularly and being able to manage their glucose levels. Therefore, as important as exercising is, showing youth and parents that being physical active is possible. That can be done in specific intensities and with proper understanding on how to handle glucose control over time.
Diabetes and youth
Youth dealing with diabetes face unique challenges. These can include the emotional, psychological, and social implications as a result of managing a chronic condition at such a young age. Both forms of diabetes can have long-term risks for complications like cardiovascular disease, kidney damage, and nerve damage. This is why early intervention and impactful management important to improving long-term health. Physical activity plays a vital role in managing diabetes, improving glycemic control, and enhancing the overall quality of life for children and adolescents.
Younger children and adolescents with T1D tend to be more physically active compared to older individuals. This possibly because of the more demanding schedules and sedentary lifestyles of older adolescents, who are often in high school or university. This age-related decrease in physical activity is a potential challenge to keeping up with the ideal healthy lifestyle. People who have had T1D for a longer amount of time also have a tendency to be less active. As diabetes progresses, people may face more barriers to engaging in physical activity. Examples of this could include anxiety about experiencing hypoglycemic events during exercise or the physical challenges posed by the long-term complications that diabetes cause. Increased physical activity in youth with T1D can be associated with improved health. These outcomes can include better lipid profiles (higher HDL-C and lower triglycerides), healthier body composition (reduced waist circumference and BMI), and improved overall physical health. These benefits are especially important during childhood and adolescence because this is when proper growth and development are occurring.
Younger people with type 2 diabetes have a tendency to have lower levels of physical activity and CRF compared to their peers without diabetes. This contributes to their poorer overall health and increases the risk of cardiovascular and metabolic complications. Despite recommendations for physical activity as part of diabetes management, many youth and young adolesents with type 2 diabetes do not meet the guidelines, hindering their ability to effectively manage blood glucose levels and improve their health. CRF is a key health indicator. Higher levels of CRF is associated with better health outcomes. This means that increasing CRF through exercise can provide important benefits for managing type 2 diabetes. There is a need for targeted interventions that promote physical activity and improve CRF in youth with type 2 diabetes to help reduce the risk of long-term complications.
When it comes to resistance training, it is found to have no significant effect on insulin sensitivity in children and adolescents, despite it having positive trends. Intervention length, training intensity, and the participants' physical maturation might explain the mixed results. Longer and higher-intensity programs showed more promising results. Future research could focus on more dire metabolic conditions like type II diabetes, investigate the role of physical maturation, and think about including longer intervention periods. While resistance training complements aerobic exercise, its standalone effects on insulin sensitivity remain unclear.
Medications
Glucose control
Most medications used to treat diabetes act by lowering blood sugar levels through different mechanisms. There is broad consensus that when people with diabetes maintain tight glucose control – keeping the glucose levels in their blood within normal ranges – they experience fewer complications, such as kidney problems or eye problems. There is, however, debate as to whether this is appropriate and cost effective for people later in life in whom the risk of hypoglycemia may be more significant.
There are a number of different classes of anti-diabetic medications. Type 1 diabetes requires treatment with insulin, ideally using a "basal bolus" regimen that most closely matches normal insulin release: long-acting insulin for the basal rate and short-acting insulin with meals. Type 2 diabetes is generally treated with medication that is taken by mouth (e.g. metformin) although some eventually require injectable treatment with insulin or GLP-1 agonists.
Metformin is generally recommended as a first-line treatment for type 2 diabetes, as there is good evidence that it decreases mortality. It works by decreasing the liver's production of glucose, and increasing the amount of glucose stored in peripheral tissue. Several other groups of drugs, mainly oral medication, may also decrease blood sugar in type 2 diabetes. These include agents that increase insulin release (sulfonylureas), agents that decrease absorption of sugar from the intestines (acarbose), agents that inhibit the enzyme dipeptidyl peptidase-4 (DPP-4) that inactivates incretins such as GLP-1 and GIP (sitagliptin), agents that make the body more sensitive to insulin (thiazolidinedione) and agents that increase the excretion of glucose in the urine (SGLT2 inhibitors). When insulin is used in type 2 diabetes, a long-acting formulation is usually added initially, while continuing oral medications.
Some severe cases of type 2 diabetes may also be treated with insulin, which is increased gradually until glucose targets are reached.
Blood pressure lowering
Cardiovascular disease is a serious complication associated with diabetes, and many international guidelines recommend blood pressure treatment targets that are lower than 140/90 mmHg for people with diabetes. However, there is only limited evidence regarding what the lower targets should be. A 2016 systematic review found potential harm to treating to targets lower than 140 mmHg, and a subsequent systematic review in 2019 found no evidence of additional benefit from blood pressure lowering to between 130 – 140mmHg, although there was an increased risk of adverse events.
2015 American Diabetes Association recommendations are that people with diabetes and albuminuria should receive an inhibitor of the renin-angiotensin system to reduce the risks of progression to end-stage renal disease, cardiovascular events, and death. There is some evidence that angiotensin converting enzyme inhibitors (ACEIs) are superior to other inhibitors of the renin-angiotensin system such as angiotensin receptor blockers (ARBs), or aliskiren in preventing cardiovascular disease. Although a more recent review found similar effects of ACEIs and ARBs on major cardiovascular and renal outcomes. There is no evidence that combining ACEIs and ARBs provides additional benefits.
Aspirin
The use of aspirin to prevent cardiovascular disease in diabetes is controversial. Aspirin is recommended by some in people at high risk of cardiovascular disease; however, routine use of aspirin has not been found to improve outcomes in uncomplicated diabetes. 2015 American Diabetes Association recommendations for aspirin use (based on expert consensus or clinical experience) are that low-dose aspirin use is reasonable in adults with diabetes who are at intermediate risk of cardiovascular disease (10-year cardiovascular disease risk, 5–10%). National guidelines for England and Wales by the National Institute for Health and Care Excellence (NICE) recommend against the use of aspirin in people with type 1 or type 2 diabetes who do not have confirmed cardiovascular disease.
Surgery
Weight loss surgery in those with obesity and type 2 diabetes is often an effective measure. Many are able to maintain normal blood sugar levels with little or no medications following surgery and long-term mortality is decreased. There is, however, a short-term mortality risk of less than 1% from the surgery. The body mass index cutoffs for when surgery is appropriate are not yet clear. It is recommended that this option be considered in those who are unable to get both their weight and blood sugar under control.
A pancreas transplant is occasionally considered for people with type 1 diabetes who have severe complications of their disease, including end stage kidney disease requiring kidney transplantation.
Diabetic peripheral neuropathy (DPN) affects 30% of all diabetes patients. When DPN is superimposed with nerve compression, DPN may be treatable with multiple nerve decompressions. The theory is that DPN predisposes peripheral nerves to compression at anatomical sites of narrowing, and that the majority of DPN symptoms are actually attributable to nerve compression, a treatable condition, rather than DPN itself. The surgery is associated with lower pain scores, higher two-point discrimination (a measure of sensory improvement), lower rate of ulcerations, fewer falls (in the case of lower extremity decompression), and fewer amputations.
Self-management and support
In countries using a general practitioner system, such as the United Kingdom, care may take place mainly outside hospitals, with hospital-based specialist care used only in case of complications, difficult blood sugar control, or research projects. In other circumstances, general practitioners and specialists share care in a team approach. Evidence has shown that social prescribing led to slight improvements in blood sugar control for people with type 2 diabetes. Home telehealth support can be an effective management technique.
The use of technology to deliver educational programs for adults with type 2 diabetes includes computer-based self-management interventions to collect for tailored responses to facilitate self-management. There is no adequate evidence to support effects on cholesterol, blood pressure, behavioral change (such as physical activity levels and dietary), depression, weight and health-related quality of life, nor in other biological, cognitive or emotional outcomes.
Epidemiology
An estimated 382 million people worldwide had diabetes in 2013 up from 108 million in 1980. Accounting for the shifting age structure of the global population, the prevalence of diabetes is 8.8% among adults, nearly double the rate of 4.7% in 1980. Type 2 makes up about 90% of the cases. Some data indicate rates are roughly equal in women and men, but male excess in diabetes has been found in many populations with higher type 2 incidence, possibly due to sex-related differences in insulin sensitivity, consequences of obesity and regional body fat deposition, and other contributing factors such as high blood pressure, tobacco smoking, and alcohol intake.
The WHO estimates that diabetes resulted in 1.5 million deaths in 2012, making it the 8th leading cause of death. However, another 2.2 million deaths worldwide were attributable to high blood glucose and the increased risks of cardiovascular disease and other associated complications (e.g. kidney failure), which often lead to premature death and are often listed as the underlying cause on death certificates rather than diabetes. For example, in 2017, the International Diabetes Federation (IDF) estimated that diabetes resulted in 4.0 million deaths worldwide, using modeling to estimate the total number of deaths that could be directly or indirectly attributed to diabetes.
Diabetes occurs throughout the world but is more common (especially type 2) in more developed countries. The greatest increase in rates has, however, been seen in low- and middle-income countries, where more than 80% of diabetic deaths occur. The fastest prevalence increase is expected to occur in Asia and Africa, where most people with diabetes will probably live in 2030. The increase in rates in developing countries follows the trend of urbanization and lifestyle changes, including increasingly sedentary lifestyles, less physically demanding work and the global nutrition transition, marked by increased intake of foods that are high energy-dense but nutrient-poor (often high in sugar and saturated fats, sometimes referred to as the "Western-style" diet). The global number of diabetes cases might increase by 48% between 2017 and 2045.
As of 2020, 38% of all US adults had prediabetes. Prediabetes is an early stage of diabetes.
History
Diabetes was one of the first diseases described, with an Egyptian manuscript from 1500 BCE mentioning "too great emptying of the urine." The Ebers papyrus includes a recommendation for a drink to take in such cases. The first described cases are believed to have been type 1 diabetes. Indian physicians around the same time identified the disease and classified it as madhumeha or "honey urine", noting the urine would attract ants.
The term "diabetes" or "to pass through" was first used in 230 BCE by the Greek Apollonius of Memphis. The disease was considered rare during the time of the Roman empire, with Galen commenting he had only seen two cases during his career. This is possibly due to the diet and lifestyle of the ancients, or because the clinical symptoms were observed during the advanced stage of the disease. Galen named the disease "diarrhea of the urine" (diarrhea urinosa).
The earliest surviving work with a detailed reference to diabetes is that of Aretaeus of Cappadocia (2nd or early 3rdcentury CE). He described the symptoms and the course of the disease, which he attributed to the moisture and coldness, reflecting the beliefs of the "Pneumatic School". He hypothesized a correlation between diabetes and other diseases, and he discussed differential diagnosis from the snakebite, which also provokes excessive thirst. His work remained unknown in the West until 1552, when the first Latin edition was published in Venice.
Two types of diabetes were identified as separate conditions for the first time by the Indian physicians Sushruta and Charaka in 400–500 CE with one type being associated with youth and another type with being overweight. Effective treatment was not developed until the early part of the 20th century when Canadians Frederick Banting and Charles Best isolated and purified insulin in 1921 and 1922. This was followed by the development of the long-acting insulin NPH in the 1940s.
Etymology
The word diabetes ( or ) comes from Latin , which in turn comes from Ancient Greek (), which literally means "a passer through; a siphon". Ancient Greek physician Aretaeus of Cappadocia (fl. 1stcentury CE) used that word, with the intended meaning "excessive discharge of urine", as the name for the disease. Ultimately, the word comes from Greek (), meaning "to pass through", which is composed of - (-), meaning "through" and (), meaning "to go". The word "diabetes" is first recorded in English, in the form diabete, in a medical text written around 1425.
The word mellitus ( or ) comes from the classical Latin word , meaning "mellite" (i.e. sweetened with honey; honey-sweet). The Latin word comes from -, which comes from , meaning "honey"; sweetness; pleasant thing, and the suffix -, whose meaning is the same as that of the English suffix "-ite". It was Thomas Willis who in 1675 added "mellitus" to the word "diabetes" as a designation for the disease, when he noticed the urine of a person with diabetes had a sweet taste (glycosuria). This sweet taste had been noticed in urine by the ancient Greeks, Chinese, Egyptians, and Indians.
Society and culture
The 1989 "St. Vincent Declaration" was the result of international efforts to improve the care accorded to those with diabetes. Doing so is important not only in terms of quality of life and life expectancy but also economicallyexpenses due to diabetes have been shown to be a major drain on healthand productivity-related resources for healthcare systems and governments.
Several countries established more and less successful national diabetes programmes to improve treatment of the disease.
Diabetes stigma
Diabetes stigma describes the negative attitudes, judgment, discrimination, or prejudice against people with diabetes. Often, the stigma stems from the idea that diabetes (particularly Type 2 diabetes) resulted from poor lifestyle and unhealthy food choices rather than other causal factors like genetics and social determinants of health. Manifestation of stigma can be seen throughout different cultures and contexts. Scenarios include diabetes statuses affecting marriage proposals, workplace-employment, and social standing in communities.
Stigma is also seen internally, as people with diabetes can also have negative beliefs about themselves. Often these cases of self-stigma are associated with higher diabetes-specific distress, lower self-efficacy, and poorer provider-patient interactions during diabetes care.
Racial and economic inequalities
Racial and ethnic minorities are disproportionately affected with higher prevalence of diabetes compared to non-minority individuals. While US adults overall have a 40% chance of developing type 2 diabetes, Hispanic/Latino adults chance is more than 50%. African Americans also are much more likely to be diagnosed with diabetes compared to White Americans. Asians have increased risk of diabetes as diabetes can develop at lower BMI due to differences in visceral fat compared to other races. For Asians, diabetes can develop at a younger age and lower body fat compared to other groups. Additionally, diabetes is highly underreported in Asian American people, as 1 in 3 cases are undiagnosed compared to the average 1 in 5 for the nation.
People with diabetes who have neuropathic symptoms such as numbness or tingling in feet or hands are twice as likely to be unemployed as those without the symptoms.
In 2010, diabetes-related emergency room (ER) visit rates in the United States were higher among people from the lowest income communities (526 per 10,000 population) than from the highest income communities (236 per 10,000 population). Approximately 9.4% of diabetes-related ER visits were for the uninsured.
Naming
The term "type 1 diabetes" has replaced several former terms, including childhood-onset diabetes, juvenile diabetes, and insulin-dependent diabetes mellitus. Likewise, the term "type 2 diabetes" has replaced several former terms, including adult-onset diabetes, obesity-related diabetes, and noninsulin-dependent diabetes mellitus. Beyond these two types, there is no agreed-upon standard nomenclature.
Diabetes mellitus is also occasionally known as "sugar diabetes" to differentiate it from diabetes insipidus.
Other animals
Diabetes can occur in mammals or reptiles. Birds do not develop diabetes because of their unusually high tolerance for elevated blood glucose levels.
In animals, diabetes is most commonly encountered in dogs and cats. Middle-aged animals are most commonly affected. Female dogs are twice as likely to be affected as males, while according to some sources, male cats are more prone than females. In both species, all breeds may be affected, but some small dog breeds are particularly likely to develop diabetes, such as Miniature Poodles.
Feline diabetes is strikingly similar to human type 2 diabetes. The Burmese, Russian Blue, Abyssinian, and Norwegian Forest cat breeds are at higher risk than other breeds. Overweight cats are also at higher risk.
The symptoms may relate to fluid loss and polyuria, but the course may also be insidious. Diabetic animals are more prone to infections. The long-term complications recognized in humans are much rarer in animals. The principles of treatment (weight loss, oral antidiabetics, subcutaneous insulin) and management of emergencies (e.g. ketoacidosis) are similar to those in humans.
| Biology and health sciences | Illness and injury | null |
33192846 | https://en.wikipedia.org/wiki/Rickshaw | Rickshaw | Rickshaw originally denoted a pulled rickshaw, which is a two- or three-wheeled cart generally pulled by one person carrying one passenger. The first known use of the term was in 1879. Over time, cycle rickshaws (also known as pedicabs or trishaws), auto rickshaws, and electric rickshaws were invented, and have replaced the original pulled rickshaws, with a few exceptions for their use in tourism.
Pulled rickshaws created a popular form of transportation, and a source of employment for male labourers, within Asian cities in the 19th century. Their appearance was related to newly acquired knowledge of ball-bearing systems. Their popularity declined as cars, trains and other forms of transportation became widely available.
Auto rickshaws are becoming more popular in some cities in the 21st century as an alternative to taxis because of their low cost of hire. Bangladesh holds the record of hosting highest number of rickshaws in the world with 40,000 rickshaws operating in the capital Dhaka alone every day. In 2023, UNESCO listed rickshaws and rickshaw art as 'intangible heritage' of Bangladesh.
Etymology
Rickshaw originates from the Japanese word jinrikisha (, jin = human, riki = power or force, sha = vehicle), which literally means "human-powered vehicle".
History
Origin
The first rickshaws were invented in France in the late 17th century, to fulfill, along with other types of carriages such as cabriolets and fiacres, the unmet demand for public transportation created by the 1679 cessation of Paris' first omnibus service. These vehicles, called "vinaigrettes" for their resemblance to the handcarts used by contemporary vinegar-sellers, were fully-enclosed two-wheeled carriages with space for a single person. Usually, they were moved by two people; one holding the bars at the front and the other pushing from behind.
A painting named "Les deux carosses" by Claude Gillot shows 2 rickshaws in 1707.
Rickshaws were independently invented in Japan circa 1869, after the lifting of a ban on wheeled vehicles from the Tokugawa period (1603–1868), and at the beginning of a period of rapid technical advancement in Japan.
Inventor
There are many theories about the inventor, with the most likely and widely accepted theory describing the rickshaw as having been invented in Japan in 1869, by Izumi Yosuke, who formed a partnership with Suzuki Tokujiro and Takayama Kosuke to build the vehicles, having been "inspired by the horse carriages that had been introduced to the streets of Tokyo a few years earlier".
Other theories about the inventor of the rickshaw include:
Reverend Jonathan Goble (sometimes called Jonathan Scobie), an American Free Baptist minister and missionary to Japan, is said to have invented the rickshaw around 1869 to transport his invalid wife through the streets of Yokohama.
An American blacksmith named Albert Tolman is said to have invented the rickshaw, or "man drawn lorry" in 1846 in Worcester, Massachusetts, for a South American bound missionary.
In New Jersey, the Burlington County Historical Society claims an 1867 invention by carriage maker James Birch, and exhibits a Birch rickshaw in its museum.
Japan historian Seidensticker wrote of the theories:
Though the origins of the rickshaw are not entirely clear, they seem to be Japanese, and of Tokyo specifically. The most widely accepted theory offers the name of three inventors, and gives 1869 as the date of invention.
Description
The vehicle had a wooden carriage that rode on "superior Western wheels" and was a dramatic improvement over earlier modes of transportation. Whereas the earlier sedan chairs required two people, the rickshaw generally only required one. More than one person was required for hilly or mountainous areas. It also provided a smoother ride for the passenger. Other forms of vehicles at the time were drawn by animals or were wheelbarrows. The vehicle also has an collapsible hood for protection from sun or rain.
The Powerhouse Museum in Sydney, has had a rickshaw in its collection for over 120 years. It was made about 1880 and is described as:
A rickshaw, or jinrikisha, is a light, two-wheeled cart consisting of a doorless, chairlike body, mounted on springs with a collapsible hood and two shafts. Finished in black lacquer-ware over timber, it was drawn by a single rickshaw runner.
Late 19th century
In the Late 19th century, hand-pulled rickshaws became an inexpensive, popular mode of transportation across Asia. Peasants who migrated to large Asian cities often worked first as a rickshaw runner. It was "the deadliest occupation in the East, [and] the most degrading for human beings to pursue."
Japan
Starting in 1870, the Tokyo government issued a permit to build and sell 人力車 (jinrikisha: rickshaw in Japanese) to the trio that are believed in Asia to be the rickshaw's inventors: Izumi Yosuke, Takayama Kosuke, and Suzuki Tokujiro. In order to operate a rickshaw in Tokyo, a seal was required from these men. By 1872, they replaced the kago and norimono, becoming the main mode of transportation in Japan, with about 40,000 rickshaws in service. At that time man-power was much cheaper than horse-power; horses were generally only used by the military. Some of the rickshaws were artistically decorated with paintings and rear elevations. In this time, the more exuberant styles of decorations were banned. If the families were well-off financially they might have their own rickshaw runner. Generally, runners covered in a day, at an average traveling speed of .
Japanese rickshaw manufacturers produced and exported rickshaws to Asian countries and South Africa.
Singapore
Singapore received its first rickshaws in 1880 and soon after they were prolific, making a "noticeable change in the traffic on Singapore's streets." Bullock carts and gharries were used before rickshaws were introduced.
Many of the poorest individuals in Singapore in the late nineteenth century were poverty-stricken, unskilled people of Chinese ancestry. Sometimes called coolies, the hardworking men found that pulling a rickshaw was a new opportunity for employment.
In 1897, martial law was declared to end a four-day rickshaw workers' strike.
Other
In China, the rickshaw was first seen in 1873 and was used for public transportation the following year. Within a year there were 10,000 rickshaws in operation. Around 1880 rickshaws appeared in India, first introduced in Simla by Reverend J. Fordyce. At the turn of the century they were introduced in Calcutta, India, and by 1914 were a conveyance for hire. The rickshaw was also introduced to Korea in the late 19th century.
20th century
After World War II, there was a major shift in the use of man-powered rickshaws:
Hand-pulled rickshaws became an embarrassment to modernizing urban elites in the Third World, and were widely banned, in part because they were symbolic, not of modernity, but of a feudal world of openly marked class distinctions. Perhaps the seated rickshaw passenger is too close to the back of the laboring driver, who, besides, is metaphorically a draught animal harnessed between shafts.
The cycle rickshaw was built in the 1880s and was first used with regularity starting in 1929 in Singapore. They were found in every south and east Asian country by 1950. By the late 1980s there were estimated 4 million cycle rickshaws in the world.
Africa
Rickshaws were introduced to Durban, South Africa, and by 1904 there were about 2,000 registered rickshaw pullers. Rickshaws operated in Nairobi in the beginning of the 20th century; pullers went on strike there in 1908. In the 1920s, they were used in Bagamoyo, Tanga, Tanzania and other areas of East Africa for short distances.
Asia
The rickshaw's popularity in Japan had declined by the 1930s with the advent of motorized forms of transportation like automobiles and trains. After World War II, when gasoline and automobiles were scarce, they made a temporary comeback. The rickshaw tradition has stayed alive in Kyoto and Tokyo's geisha districts. In the 1990s, German-made cycle rickshaws called "velotaxis" were introduced in Japanese cities, including Kobe.
In post-war Hong Kong, rickshaws was one of the main transportation either for transporting goods or for transporting people during the Japanese invasion, known as the Battle of Hong Kong. Japanese military hired many rickshaw pullers to have them gathered and organize with other cooks and seamen for an underground armed team to enact the anti-British Colony clan. However, after World War II, other forms of transport such as pedicabs and streetcar became strong competitors of rickshaws, leading the business of rickshaws into stagnation.
Chiuchow men formed a faction within the Canton rickshaw pullers union in the 1920s. In addition, the Chiuchow pullers could be identified by their hats which were rounded and flat at the top; while the rival Hoklo men had a cone-shaped headgear with sharp points at the top. However, rickshaw use began to decline in the 1920s as the government introduced the streetcar system in 1924. The number of rickshaw pullers had declined from 44,200 to 25,877 six months after the opening of the tramway. It had also caused the Beijing tramway riot in October 1929. A rough form of a rickshaw is sometimes used for hauling coal, building materials or other material. Both motorized and pedal-power cycle rickshaws, or pedicabs, were used for short-distance passenger travel. There are still many rickshaws in many cities for either touring purposes (in big cities such as Beijing and Shanghai, with traditional Chinese rickshaws) or short range transportation in some counties. However, the new Communist government banned rickshaws in Canton in the early 1950s, leading to another low tide of rickshaws.
In Singapore, the rickshaw's popularity increased into the 20th century. There were approximately 50,000 rickshaws in 1920 and that number had doubled by 1930. Cycle rickshaws were used in Singapore beginning in 1929. Within six years pulled rickshaws were outnumbered by cycle rickshaws, which were also used by sightseeing tourists.
In the 1930s, cycle rickshaws were used in Dhaka, Bangladesh, Kolkata, India; and Jakarta, Indonesia. By 1950 they could be found in many South and East Asian countries. By the end of the 20th century, there were 300,000 such vehicles in Dhaka. By the end of 2013, there were about 100,000 electric rickshaws in Delhi.
In Viêtnam, the rickshaw was called with the French name "Pousse-pousse". In 1883, Jean Thomas Raoul Bonnal, Supérior Résident of Tonkin, import 2 rickshaws from Japan in the city of Hanoi and made copies of them.
North America
Pedicabs were introduced in North America in 1962, where they were a means of transportation at the Seattle World's Fair in the state of Washington.
21st century
The 21st century has seen a resurgence in rickshaws, particularly in motorized rickshaws and cycle rickshaws. Auto rickshaws, also called velotaxis, have resurged as they are about 1/3 to 1/2 the cost of regular taxis. German velotaxis are three-wheeled, powered vehicles with a space for a driver and, behind the driver, space for two passengers. Cycle rickshaws are used in many Asian, North American, and European cities. They are increasingly being used as an eco-friendly way of short-range transportation, particularly in urban areas. Along with auto rickshaws, they are also used (particularly by Asian cities) for tourism, because of their "novelty value as an entertaining form of transportation".
Africa
In Madagascar, pulled cycle and auto rickshaws are a common form of transportation in a number of cities, especially Antsirabe. They are known as pousse-pousse, meaning push-push. Aboboyaa is a tricycle used in the transportation of goods and service in Ghana.
Asia
Macau still uses tri-wheeled bicycle rickshaw, or riquexó in Portuguese, as Macau was a Portuguese colony in the past. This kind of transportation was very famous until the late 20th century, due to the fact of being a small city and few cars, not so many motorcycles, very bad public transport and no other transport such as train or subway. You can go around Macau peninsula and the twos island on rickshaw, and visit the Riquexó Museum and see the evolution of rickshaws from the 18th century until modern times.
Automated cycle rickshaws, called velotaxis, are popular in Kyoto and Tokyo, Japan. Their use is growing at a rate of about 20–30% a year in Japanese cities. The traditional rickshaws are still alive for travelers in some tourist places in Japan. Rickshaws are found in Hong Kong. In China, automated and pedal-power cycle rickshaws, or pedicabs, are used for short-distance passenger travel in large cities and many medium-sized cities. Most Indian cities offer auto rickshaw service; hand-pulled rickshaws do exist in some areas, such as Kolkata (Calcutta) as a part of their transport system which also includes cycle rickshaws. Sri Lanka has over 1 million auto rickshaws registered in use as of 2018.
Australia
In Australia, cycle rickshaws or trishaws (3 wheels) are used in Melbourne and St Kilda. They are also seen in Cowaramup, Western Australia at Bakehouse '38.
Europe
Cycle rickshaws or trishaws (3 wheels) are used in some large continental European cities, such as:
Austria: Vienna
Denmark: Copenhagen and Odense
Estonia: Tallinn
France: Paris and Nantes
Germany: Lake Constance, Berlin, Frankfurt, Hamburg, Hanover and other cities.
Hungary: Budapest
Ireland: Cork and Dublin.
Italy: Florence, Milan, and Rome
The Netherlands: Amsterdam and in the Caribbean, Willemstad
Norway: Oslo
Poland: Kraków and Łódź
Russia: Saint Petersburg
Spain: Barcelona
Switzerland: Bern
Romania: Bucharest
Within the United Kingdom, pedicabs operate in:
London, mostly in Soho and other areas of central London. Their activity is not regulated by Transport for London, so there are no limits of the fees they can charge. The Queen's Speech in May 2022 introduced plans to legislate a licensing scheme for pedicabs.
Edinburgh, where vendors are hired like taxis and provide tours.
Oxford.
North America
In Canada, there are pedicabs in operation in Victoria and Vancouver. They are regulated in Toronto and Vancouver. Pulled rickshaw rides are available in downtown Ottawa, with tours of historical Byward Market.
In the United States, San Diego and New York City each host hundreds of pedicabs; dozens of other North American cities also have pedicab services. In New York, human powered transport is used primarily by tourists due to its cost. In New Orleans, pedicabs have been used to transport French Quarter tourists since the summer of 2012.
In Mexico, there are thousands of pedicabs. All drivers are in informal circumstances, and have precarious working conditions, long hours (11.3 hours a day), low wages (US$59.18 per week), and no social protections or benefits. 6.3% reported suffering from a disease, 49.5% corresponded to musculoskeletal conditions and only 11.6% were affiliated to any health system. 53.8% are owners of the vehicle and, although it does not seem to influence physical illness (P=0.03), it is related to the psychosocial ones (P=0.260).
Types
Types of rickshaws include:
a pulled rickshaw; a two-wheeled passenger cart pulled by a human runner
a cycle rickshaw, also called a pedicab
an auto rickshaw, also called a tuk-tuk, auto, mototaxi, or baby taxi
an electric rickshaw, also called e-rickshaw.
| Technology | Road transport | null |
25970423 | https://en.wikipedia.org/wiki/IPad | IPad | The iPad is a brand of iOS- and iPadOS-based tablet computers that are developed and marketed by Apple. The first-generation iPad was introduced on January 27, 2010. Since then, the iPad product line has been expanded to include the smaller iPad Mini, the lighter and thinner iPad Air, and the flagship iPad Pro models. As of 2022, over 670 million iPads have been sold, making Apple the largest vendor of tablet computers. Due to its popularity, the term "iPad" is sometimes used as a generic name for tablet computers.
The iPhone's iOS operating system (OS) was initially used for the iPad, but in September 2019, its OS was switched to a fork of iOS called iPadOS that has better support for the device's hardware and a user interface tailored to the tablets' larger screens. Since then, major versions of iPadOS have been released annually. The iPad's App Store is subject to application and content approval. Many older devices are susceptible to jailbreaking, which circumvents these restrictions.
The original iPad was well-received for its software and was recognized as one of the most-influential inventions of 2010. As of the third quarter of 2021, the iPad had a market share of 34.6% among tablets. Beside personal use, the iPad is used in the business, education, healthcare, and technology sectors. There are two connectivity variants of iPad; one has only Wi-Fi, and one has additional support for cellular networks. Accessories for the iPad include the Apple Pencil, Smart Case, Smart Keyboard, Smart Keyboard Folio, Magic Keyboard, and several adapters.
History
Background
Apple co-founder and CEO Steve Jobs said in a 1983 speech: "What we want to do is we want to put an incredibly great computer in a book that you can carry around with you and learn how to use in 20 minutes... and we really want to do it with a radio link in it so you don't have to hook up to anything and you're in communication with all of these larger databases and other computers".
In 1993, Apple worked on the Newton MessagePad, a tablet-like personal digital assistant (PDA). John Sculley, Apple's chief executive officer, led the development. The MessagePad was poorly received for its indecipherable handwriting recognition feature and was discontinued at the direction of Jobs, who returned to Apple in 1998 after an internal power struggle. Apple also prototyped a PowerBook Duo–based tablet computer but decided not to release it to avoid hurting MessagePad sales.
In May 2004, Apple filed a design trademark patent in Europe for a handheld computer, hypothetically referencing the iPad, beginning a new round of speculation that led to a 2003 report of Apple-affiliated manufacturer Quanta leaking Apple's orders for wireless displays. In May 2005, Apple filed US Design Patent No. D504,889 that included an illustration depicting a man touching and using a tablet device. In August 2008, Apple filed a 50-page patent application that includes an illustration of hands touching and gesturing on a tablet computer. In September 2009, Taiwan Economic News, citing "industry sources", reported the tablet computer Apple was working on would be announced in February 2010, although the announcement was made in that year's January.
The iPad's concept predates that of the iPhone, although the iPhone was developed and released before the iPad. In 1991, Apple's chief design officer Jonathan Ive devised an industrial design of a stylus-based tablet, the Macintosh Folio, which led to the development of a larger tablet prototype project codenamed K48 that Apple began in 2004. Ive sought to develop the tablet first but came to an agreement with Jobs the iPhone was more important and should be prioritized.
iPad
The first generation of iPad was announced on January 27, 2010, and pre-ordering began on March 12. Initial reaction to the product name, which struck some women as a menstrual pun "indicative of a male-helmed team oblivious" to the connotations, was negative.
A Wi-Fi-capable version was released in the United States on April 3 and a 3G-capable version was released on April 30. Apple released iPad models internationally on May 28, July 23, and September 17. The first iPad has a 1 GHz Apple A4 CPU with 256 MB of RAM and a PowerVR SGX535 GPU. It has four buttons; a home button that directs the user to its homepage, a wake-and-sleep button, and two volume-control buttons. Its multi-touch-based display has a resolution of 1,024 by 768 pixels.
The second generation of iPad was announced on March 2, 2011, and released on March 11. It is 33% thinner and 15% lighter than its predecessor, and uses a dual-core Apple A5 chip consisting of a twice-as-fast CPU and a nine-times-faster GPU. It has one camera each on the front and the back, both of which support Apple's video-telephony service, FaceTime. Apple slimmed the iPad by eliminating the display's stamped-sheet-metal frame, using thinner glass for the screen overlay, and eliminating some space between the display and battery.
The third generation of iPad was announced on March 7, 2012, and released on March 16. It uses a dual-core Apple A5X chip embedded with quad-core graphics. Its Retina Display is 2,048 by 1,536 pixels and its pixels are 50% denser than those of standard displays. Unlike the iPhone and iPod Touch's built-in applications, which work in portrait, landscape-left and landscape-right orientations, the iPad's built-in applications support the upside-down orientation of the device. Consequently, the device has no "native" orientation; only the relative position of the home button changes.
The fourth generation of iPad was announced on October 23, 2012, and released on November 2. It has an Apple A6X chip, improved LTE and WiFi connectivity, a five-megapixel, rear-facing camera that is capable of recording 1080p videos, and a 720p front-facing FaceTime HD camera. Its display has a resolution of 2,048 by 1,536 pixels.
The fifth generation of iPad was announced on March 21, 2017, and released on March 24. It uses an Apple A9 chip with an accompanying M9 motion coprocessor, and its cameras can capture low-light and HD-quality shots. Despite using the same Apple A9 and M9 processors as the 2015 iPhone 6S, it lacks support of the always-on "Hey Siri" voice recognition, a feature advertised as being made possible by low-power processing in those chips.
The sixth generation of iPad was announced and released on March 27, 2018. It uses a dual-core Apple A10 Fusion chip and feature a 1080p and 30fps rear mounted 8-megapixel iSight camera, and a 720p Facetime HD camera. It was the first non-Pro iPad to support the Apple Pencil. It also had faster FaceTime HD, LTE connectivity, Touch ID, and multitask functionalities.
The seventh generation of iPad was announced on September 10, 2019, and released on September 25. It uses a 64-bit Apple A10 Fusion chip, with a 4-core CPU and 6-core GPU. Its slightly larger 10.2" Retina Display has a resolution of 2,160 × 1,620 (3.5 million pixels). It added support for the Smart Keyboard accessory.
The eighth generation of iPad was announced on September 15, 2020, and released on September 18. It uses an Apple A12 Bionic chip, which has a 40 percent faster 6-core CPU and a 2-times faster 4-core GPU than its predecessor. The Apple A12 also included an embedded Neural Engine, and is capable of processing 5 trillion operations per second.
The ninth generation of iPad was announced and released on September 14, 2021. It uses an Apple A13 Bionic chip, which has a 20% faster CPU and GPU and an embedded, artificial intelligence–immersed Neural Engine. Its 12-megapixel ultra wide front camera adds support for Apple's "Center Stage Mode" technology, which locates people in the frame and tracks the camera view to keep them centered. Its Retina Display added support for the True Tone technology, which automatically adjusts the screen color temperature according to the ambient lighting.
The 10th-generation iPad was announced on October 18, 2022, with pre-orders starting immediately and availability set for October 26. It uses the Apple A14 Bionic chip, has a larger 10.9-inch screen, and replaces the Lightning connector with USB-C. Unlike all previous models in the iPad range, as well as the sixth-generation iPad Pro announced the same day, this model's front-facing camera is placed along the device's long edge, making it more suitable for video calling applications. Despite having a USB-C connector, it is not compatible with the second-generation Apple Pencil that can be used with all other USB-C iPads, instead using the first-generation Pencil with a USB-C-to-Lightning adapter, which will be included with new Pencil purchases. While lacking the Smart Connector of the Pro and Air lines, it is compatible with a new Magic Keyboard Folio announced alongside the device. This model did not immediately replace the 9th-generation iPad; Apple continued to sell the older model at the same price, while the price for the newer 10th-generation model was increased.
iPad Mini
The first generation of the flagship, smaller iPad Mini was announced on October 23, 2012, and released on November 2. The 7.9-inch display has a resolution of 1024 by 768 pixels. The tablet uses a dual-core Apple A5 chip, and has hardware resembling that of the second generation of iPad. It has a FaceTime HD camera, a 5-megapixel iSight camera, an ultrafast wireless LTE range, and a 802.11a/b/g/n standard Wifi connectivity. It targets the emerging sector of mini tablets, such as Kindle Fire and Nexus 7.
The second generation of iPad Mini was announced on October 22, 2013, and released on November 12. Its hardware resembles that of the first generation of iPad Air.
The third generation of iPad Mini was announced on October 16, 2014, and released on October 22. It uses an Apple A7 chip with an embedded M7 motion coprocessor, and its 7.9-inch Retina screen display has a resolution of 2048 by 1536 pixels. It includes a 1080p HD camera, a FaceTime HD camera, and a 5-megapixel iSight camera.
The fourth generation of iPad Mini was announced and released on September 9, 2015. It uses a dual-core Apple A8 chip with an embedded Apple M8 motion coprocessor. Its headphone jack was re-positioned with the removal of a mute switch.
The fifth generation of iPad Mini was announced and released on March 18, 2019. It uses an Apple A12 Bionic chip, with a 3-times faster CPU and a 9-times faster GPU than its predecessor. It features a Truetone-based Retina screen display with 25% wider Color and higher pixel density.
The sixth generation of iPad Mini was announced and released on September 24, 2021. The display size was increased to 8.3 inches. It uses an Apple A15 Bionic chip, with a 40% faster 6-core CPU and 80% faster 5-core GPU. Its 16-core Neural Engine and AI accelerators within the CPU delivers a 2× boost of AI performance. Its 12-megapixel Ultra Wide front camera featured Apple's "Center Stage Mode" technology, while its 12-megapixel back camera had larger apertures, True Tone flash, and Smart HDR automatic shadow and highlight recovery. It includes a USB-C port, capable of transferring up to 5 gigabits per second of data; improved landscape stereo speakers; and a brighter Liquid Retina Display.
The seventh generation of iPad Mini was announced on October 15, 2024. It is powered by an Apple A17 Pro chip, which Apple says has 30 percent faster CPU, 25 percent faster GPU and a neural engine twice as fast as the previous generation. This model is marketed by Apple as built for Apple Intelligence, its upcoming artificial intelligence features. The storage on the lowest priced variant has increased from 64GB of previous generation to 128GB. The new Wi-Fi 6E chip is faster, the USB-C port is faster compared to previous generations. Like the iPad Pro (7th generation) and the iPad Air (6th generation), it is only compatible with the Apple Pencil Pro alongside the lower-entry Apple Pencil with USB-C port, making it incompatible with the Apple Pencil (2nd generation). It is not compatible with the Magic Keyboard for iPad and Smart Keyboard Folio due to its smaller form factor.
iPad Air
The first generation of iPad Air was announced on October 22, 2013, and released on November 1. It used an Apple A7 chip with an embedded Apple M7 Motion coprocessor; the chip included over a billion transistors and comprised a 2× faster CPU and GPU. It debuted the 802.11n-based MINO technology used in its Wi-Fi connectivity, and it had an extended range of LTE telecommunication. It also came with a Retina Display.
The second generation of iPad Air was announced on October 16, 2014, and released on October 22. It used an Apple A8X chip with a 2.5× faster CPU. Its 8MP iSight Camera had 1.12-micron pixels and a f/2.4 aperture, while its FaceTime Camera had f/2.2 aperture and 81% light capacity. Its display had a revised 56% lower reflective rate. It also had an extended range of LTE telecommunication service.
The third generation of iPad Air was announced on March 18, 2019, and released on March 25. It used an Apple A12 Bionic with an embedded Neural Engine, 6-core CPU and 4-core GPU. Its 866 Mbit/s WiFi connectivity are LTE-based, and it is equipped with a 1080p HD video camera.
The fourth generation of iPad Air was announced on September 15, 2020, and released on October 23. It used an Apple A14 Bionic chip, which comprised 11.9 billion transistors, a 40% faster 6-core CPU, a 30% faster 4-core GPU, and an embedded Neural Engine that can process 11 trillion operations per second. Its 10.9-inch Liquid Retina Screen display has a resolution of 2360 by 1640 pixels (3.8 million pixels). Its front 7-megapixel Facetime Camera is of 1080p and 60 fps, while its 12-megapixel webcam featured 8 aperture, 4K, 60fps, and video stabilization. The fourth generation iPad Air replaces the Lightning connector with USB-C.
The fifth generation of iPad Air was announced on March 8, 2022, and released on March 18. It used an Apple M1 chip.
The sixth generation of iPad Air was announced on May 7, 2024, with the general availability on May 15, 2024. It used an Apple M2 chip. It is the first iPad Air to have the available two display size options including 11-inch and 13-inch. It features the Apple Pencil hover which is previously featured on the iPad Pro (6th generation) and the landscape oriented front-facing camera which is previously used on the tenth-generation iPad. Like the iPad Pro (7th generation), it is only compatible with the Apple Pencil Pro alongside the lower-entry Apple Pencil with USB-C port, making it incompatible with the Apple Pencil (2nd generation). Unlike the iPad Pro (7th generation), it is only compatible with the Magic Keyboard that was designed for the iPad Pro 3rd, 4th, 5th and 6th generations.
iPad Pro
The first generation of the high-end and professional flagship iPad Pro was announced on September 9, 2015, and released on November 11, (12.9-inch version) and March 31 (9.7-inch). It used an Apple A9X chip, with a 2× higher memory bandwidth and a 1.8× faster CPU than its predecessor. Its audio system consisted of 4 audio ports and its volume were more 3× more efficient than the second generation of iPad Air, and its 12-inch screen display had a resolution of 2732 by 2043 pixels.
The second generation of iPad Pro was announced on June 5, 2017, and was released on June 13. It used an Apple A10X chip, with a 6-core CPU and 12-core GPU. It can process 120 Hz HDR quality medias, 2× higher quality than its predecessor. Its ultra-low reflective Retina Display featured a 50% optimized True Tone technology (which automatically adjust the screen accordingly to its ambient color and brightness rates), Wide Color Integration, and up to 500 nit brightness rates. It also had a 12-megapixel rear-facing camera and a 7-megapixel front-facing camera.
The third generation of iPad Pro was announced on October 30, 2018, and released on November 7, and it is the first iPad to support 1TB of storage. It used a 7 nm Apple A12X Bionic chip, which comprised 11 billion transistors, an 8-core CPU, 7-core GPU and an embedded Neural Engine capable of processing 5 trillion operations per second. Apple replaced the Touch ID fingerprint recognition biometric authentication with its facial counterpart, Face ID.
The fourth generation of iPad Pro was announced and released on March 18, 2020. It used an Apple A12Z chip, with an 8-core CPU and 8-core GPU. Its Gbit-class Wi-Fi connectivity is 60% faster than that of its predecessor. It introduced a 10-megapixel ultra-wide camera, alongside its 12-megapixel wide camera, capable of capturing 4k video. These cameras allow it to capture medias with wider visibility, and its audio system automatically detects and attracts any orientation nearby.
The fifth generation of iPad Pro was announced on April 20, 2021, and released on May 21. It used an innovative desktop-class Apple M1 chip, which comprised a 40% faster 8-core CPU, a 4× faster 8-core GPU, and a 4× higher bandwidth. It featured a ƒ/1.8 aperture 12-megapixel wide-angle pro camera (captures high quality shots) and a ƒ/2.4 aperture 10-megapixel ultrawide camera (captures enhanced Augmented Reality interactive experience). It debuted Apple's "Center Stage mode" technology, which pinpoints the positions of the users and automatically tracks the camera view accordingly to perspectivally centralize them. The 12.9-inch version had a mini LED-based Liquid Retina XDR display, compared to the 11-inch model's lesser IPS LCD-based Liquid Retina display.
The sixth generation of iPad Pro was announced on October 18, 2022, and released on October 26. It used an Apple M2 chip, with an 8-core CPU and 10-core GPU.
The seventh generation of iPad Pro was announced on May 7, 2024, a release on May 15, 2024. It uses an Apple M4, making it the first Apple device to use this chip. It is the first iPad Pro to use OLED display, called Ultra XDR Retina Display. The chassis design is about 5.3mm thin on the 11-inch model and 5.1mm thin on the 13-inch model, thinner than previous models. It features a Tandem OLED display, an updated Pro Cameras (which lacks the ultra-wide camera lens), and a landscape oriented front-facing camera, previously used in the iPad (10th generation). It is only compatible with the Apple Pencil Pro, the lower-entry Apple Pencil with USB-C port and the new, thinner Magic Keyboard, making it incompatible with the Apple Pencil (2nd generation) and the Magic Keyboard that was designed for the 3rd through 6th generations.
Hardware
Overview
Cellular connectivity
The iPad comes in two variants: Wi-Fi only and Wi-Fi with cellular support. Unlike the iPhone, the cellular variant did not support voice calls and text messages, but only data connectivity; it also had an additional micro-SIM circuit slot attached on the side. The 3G-based iPad is compatible with any GSM carrier, unlike the iPhone which is usually sold 'locked' to specific carriers. For the first generation of iPad, cellular access from T-Mobile was limited to slower EDGE cellular speeds because T-Mobile's network at the time used different frequencies.
The second generation of iPad introduced a third tier of CDMA support from Verizon, which is available separately from the AT&T-based version. The fifth generation of iPad used a nano-SIM circuit slot, while its predecessors used micro-SIM. The iPads used two frequency bands; both support the same quad-band GSM and quad-band UMTS frequencies. One supports LTE bands 4 and 17 (principally intended for use on the U.S. AT&T network), and the other supports LTE bands 1, 3, 5, 13, 25 and CDMA EV-DO Rev. A and Rev. B.
Apple extended the range of cellular compatibilities worldwide with the release of the fifth generation of iPad and the second generation of iPad Mini, worldwide and all major carriers across North America. The iPad Air and iPad Mini come in two cellular sub-variants, all of which featured nano-SIMs, quad-band GSM, penta-band UMTS, and dual-band CDMA EV-DO Rev. A and B. One supports LTE bands 1, 2, 3, 4, 5, 7, 8, 13, 17, 18, 19, 20, 25 and 26, and the other supports LTE bands 1, 2, 3, 5, 7, 8, 18, 19, 20 and TD-LTE bands 38, 39 and 40.
Accessories
Apple offers many accessories for its iPad models, ranging from keyboards, styluses, cases, to adapters; a 10 W power adapter is bundled with the device. In addition to a camera connection kit which consists of two adapters for the iPad's dock connector, one of USB Type A and one of SD card reader; these adapters can transfer photographs and videos and connect USB audio card and MIDI keyboard.
Apple's list of accessories included the Apple Pencil ― a wireless stylus pen, Smart Cover ― a magnetic screen protector that align to the face of an iPad with three folds that is convertible into a stand, Smart Case ― a fine case combining the functions of a Smart Cover and a back-protection case, Smart Keyboard Folio ― an externally-paired keyboard and a combination of a Smart Case and its predecessor, a Smart Keyboard, Magic Keyboard, ― an externally-paired keyboard similar to the formers but with integrated trackpads which the Smart Keyboard Folio and Smart Keyboard lack.
Software
Since its introduction in 2010, the iPad runs on the iPhone's iOS mobile operating system, but it was later replaced with an optimized derivation, iPadOS, in September 2019. It shares the former's development environment and many of its applications and features. The iPad is compatible with nearly every iPhone application through iOS, and developers can optimize these applications to take full advantage of the iPad's software. They used iOS SDK, a software development kit.
The iOS user interface is based upon direct manipulation, using multi-touch gestures such as swipe, tap, pinch, and reverse pinch. Interface control elements include sliders, switches, and buttons. Internal accelerometers are used by some applications to respond to shaking the device (one common result is the undo command) or rotating it in three dimensions (one common result is switching between portrait and landscape mode). Various accessibility described in § Accessibility functions enable users with vision and hearing disabilities to properly use iOS.
iOS devices boot to the homescreen, the primary navigation and information "hub" on iOS devices, analogous to the desktop found on personal computers. iOS homescreens are typically made up of app icons and widgets; app icons launch the associated app, whereas widgets display live, auto-updating content, such as a weather forecast, the user's email inbox, or a news ticker directly on the homescreen.
Along the top of the screen is a status bar, showing information about the device and its connectivity. The status bar itself contains two elements, the Control Center and the Notification Center.
iOS' Control Center can be "pulled" down from the top right of the notch, giving access to various toggles to manage the device more quickly without having to open the Settings. It is possible to manage brightness, volume, wireless connections, music player, etc. A homescreen may be made up of several pages, between which the user can swipe back and forth, one of the ways to do this is to hold down on the "dots" shown on each page and swipe left or right. To the right of the last page, the App Library lists and categorizes apps installed on the device. Apps within each category are arranged based on the frequency of their usage. In addition to a category for suggested apps, a "recent" category lists apps recently installed alongside App Clips recently accessed. Users can search for the app they want or browse them in alphabetical order.
iOS' multitasking API included Background audio – application continues to run in the background as long as it is playing audio or video content, voice over IP – application is suspended when a phone call is not in progress, Push notification, Local notifications – application schedules local notifications to be delivered at a predetermined time, Task completion – application asks the system for extra time to complete a given task, Fast app switching – application does not execute any code and may be removed from memory at any time, Newsstand – applications can download content in the background to be ready for the uses, External Accessory – application communicates with an external accessory and shares data at regular intervals, Bluetooth Accessory – application communicates with a Bluetooth accessory and shares data at regular intervals, and Background application update.
iPadOS features a multitasking system developed with more capabilities compared to iOS, with features like Slide Over and Split View that make it possible to use multiple different applications simultaneously. Double-clicking the Home Button or swiping up from the bottom of the screen and pausing will display all currently active spaces. Each space can feature a single app, or a Split View featuring two apps. The user can also swipe left or right on the Home Indicator to go between spaces at any time, or swipe left/right with four fingers.
In iPadOS, while using an app, swiping up slightly from the bottom edge of the screen will summon the Dock, where apps stored within can be dragged to different areas of the current space to be opened in either Split View or Slide Over. Dragging an app to the left or right edge of the screen will create a Split View, which will allow both apps to be used side by side. The size of the two apps in Split View can be adjusted by dragging a pill-shaped icon in the center of the vertical divider and dragging the divider all the way to one side of the screen closes the respective app. If the user drags an app from the dock over the current app, it will create a floating window called Slide Over which can be dragged to either the left or right side of the screen. A Slide Over window can be hidden by swiping it off the right side of the screen, and swiping left from the right edge of the screen will restore it. Slide Over apps can also be cycled between by swiping left or right on the Home Indicator in the Slide Over window and pulling up on it will open an app switcher for Slide Over windows. A pill-shaped icon at the top of apps in Split View or Slide Over allows them to be switched in an out of Split View and Slide Over.
Restrictions
Digital rights management
The iPad does not employ digital rights management (DRM), but the OS prevents users from copying or transferring certain content outside of Apple's platform without authorization, such as TV shows, movies, and apps. Also, the iPad's development model requires anyone creating an app for the iPad to sign a non-disclosure agreement and pay for a developer subscription. Critics argue Apple's centralized app approval process and control of the platform itself could stifle software innovation. Of particular concern to digital rights advocates is Apple's ability to remotely disable or delete apps on any iPad at any time.
Digital rights advocates, including the Free Software Foundation, Electronic Frontier Foundation, and computer engineer and activist Brewster Kahle, have criticized the iPad for its digital rights restrictions.
In an article published by NPR in April 2010, Laura Sydell, concluded that "as more consumers have fears about security on the Internet, viruses, and malware, they may be happy to opt for Apple's gated community."
In 2014, the Russian government switched from iPads to Android devices over security concerns.
Jailbreaking
Like other iOS devices, the iPad can be "jailbroken", depending on which version of iOS or iPadOS it is running, thus allowing applications and programs that are not authorized by Apple to run on the device. Once it is jailbroken, users are able to download many applications previously unavailable through the App Store via unofficial installers such as Cydia, as well as illegally pirated applications. Apple claims jailbreaking "can" void the factory warranty on the device in the United States even though jailbreaking is legal. The iPad, released in April 2010, was first jailbroken in May 2010 with the Spirit jailbreak for iOS version 3.1.2. The iPad can be jailbroken on iOS versions 4.3 through 4.3.3 with the web-based tool JailbreakMe 3.0 (released in July 2011), and on iOS versions including 5.0 and 5.0.1 using redsn0w. Absinthe 2.0 was released on May 25, 2012, as the first jailbreak method for all iOS 5.1.1 devices except the 32 nm version of the iPad 2.
Censorship
Apple's App Store, which provides iPhone and iPad applications, imposes censorship of content, which has become an issue for book publishers and magazines seeking to use the platform. The British newspaper The Guardian described the role of Apple as analogous to that of British magazine distributor WH Smith, which for many years imposed content restrictions.
Due to the exclusion of pornography from the App Store, YouPorn and others changed their video format from Flash to H.264 and HTML5 specifically for the iPad. In an e-mail exchange with Ryan Tate from Valleywag, Steve Jobs claimed that the iPad offers "freedom from porn", leading to many upset replies including Adbustings in Berlin by artist Johannes P. Osterhoff and in San Francisco during WWDC10.
Original reception
Media reaction to the original iPad was mixed. The media noted the positive response from fans of the device, with thousands of people queued on the first day of sale in a number of these countries.
The iPad was quickly successful and sold in large numbers after its 2010 launch. Analysts have noted that while Apple's previous iPod and iPhone launches took some time till taking off, the iPad was commercially popular from the beginning and faced little market competition during its first year.
Reaction to the announcement
Media reaction to the announcement of the original iPad was mixed. Walter Mossberg wrote, "It's about the software, stupid", meaning hardware features and build are less important to the iPad's success than software and user interface, his first impressions of which were largely positive. Mossberg also called the price "modest" for a device of its capabilities, and praised the ten-hour battery life. Others, including PC Advisor and the Sydney Morning Herald, wrote that the iPad would also compete with proliferating netbooks, most of which use Microsoft Windows. The base model's $499 price was lower than pre-release estimates by the tech press, Wall Street analysts, and Apple's competitors, all of whom were expecting a much higher entry price point.
CNET also criticized the iPad for its apparent lack of wireless sync which other portable devices such as Microsoft's Zune have had for a number of years. The built-in iTunes app is able to download from the Internet as well.
Critical response
Reviews of the original iPad have been generally favorable. Walt Mossberg, then of The Wall Street Journal, called it a "pretty close" laptop killer. David Pogue of The New York Times wrote a "dual" review, one part for technology-minded people, and the other part for non-technology-minded people. In the former section, he notes that a laptop offers more features for a cheaper price than the iPad. In his review for the latter audience, however, he claims that if his readers like the concept of the device and can understand what its intended uses are, then they will enjoy using the device. PC Magazine'''s Tim Gideon wrote, "you have yourself a winner" that "will undoubtedly be a driving force in shaping the emerging tablet landscape." Michael Arrington of TechCrunch said, "the iPad beats even my most optimistic expectations. This is a new category of device. But it also will replace laptops for many people." PC World criticized the iPad's file sharing and printing abilities, and Ars Technica said sharing files with a computer is "one of our least favorite parts of the iPad experience."
The media also praised the quantity of applications, as well as the bookstore and other media applications. In contrast they criticized the iPad for being a closed system and mentioned that the iPad faces competition from Android-based tablets, that outsold iPads in 2013, surpassing iPads in the second quarter of 2013, and have overtaken iPad's installed base, and has lost majority of web browsing to Android, by StatCounter estimates, in South America, Africa, most of Asiamany large countries there and in Eastern Europe. The Independent criticized the iPad for not being as readable in bright light as paper but praised it for being able to store large quantities of books. After its UK release, The Daily Telegraph said the iPad's lack of Adobe Flash support was "annoying."
Recognition
The original iPad was selected by Time magazine as one of the 50 Best Inventions of the Year 2010, while Popular Science chose it as the top gadget behind the overall "Best of What's New 2010" winner Groasis Waterboxx.
Usage
Market share
The iPad had a relatively stable tablet global market share. It received a significant drop in the third quarter of 2012 but gradually recovered, although not as abundant as before. As of the third quarter of 2021, it had a market share of 34.6%.
Business
While the iPad is mostly used by consumers, it also has been taken up by business users. Within 90 days of its release, the iPad managed to penetrate 50% of Fortune 100 companies. Some companies are adopting iPads in their business offices by distributing or making available iPads to employees. Examples of uses in the workplace include attorneys responding to clients, medical professionals accessing health records during patient exams, and managers approving employee requests.
A survey by Frost & Sullivan shows that iPad usage in office workplaces is linked to the goals of increased employee productivity, reduced paperwork, and increased revenue. The research firm estimates that "The mobile-office application market in North America may reach $6.85 billion in 2015, up from an estimated $1.76 billion [in 2010]."
Since March 2011, the US Federal Aviation Administration (FAA) has approved the iPad for in-cockpit use to cut down on the paper consumption in several airlines. In 2011, Alaska Airlines became the first airline to replace pilots' paper manuals with iPads, weighing compared to for the printed flight manuals. It hopes to have fewer back and muscle injuries. More than a dozen airlines have followed suit, including United, which has distributed iPads to cockpits. Also, many airlines now offer their inflight magazine as a downloadable application for the iPad.
Education and healthcare
The iPad has several uses in the classroom, and has been praised as a valuable tool for homeschooling and distance education. Soon after the iPad was released, it was reported that 81% of the top book apps were for children. The iPad has also been called a revolutionary tool to help children with autism learn how to communicate and socialize more easily.
In the healthcare field, iPads and iPhones have been used to help hospitals manage their supply chain. For example, Novation, a healthcare contracting services company, developed VHA PriceLynx (based on the mobile application platform of business intelligence software vendor MicroStrategy), a business intelligence app to help health care organizations manage its purchasing procedures more efficiently and save money for hospitals. Guillermo Ramas of Novation states, "Doctors won't walk around a hospital with a laptop. With an iPad it's perfect to walk around the hospital with as long as they have the information they need."
In 2013, Gianna Chien (aged 14) presented to more than 8,000 doctors at the Heart Rhythm Society meeting that the Apple iPad 2 can, in some cases, interfere with life-saving heart devices (pacemakers) because of the magnets inside. The iPad User Guide advised pacemaker users to keep iPads at least away from the pacemaker. A study in 2014 found that the iPad 2 could cause electromagnetic interference (EMI) in implantable cardioverter defibrillators.
Consumer usage
In the United States, fans attending Super Bowl XLV, the first Super Bowl since the iPad was released, could use an official National Football League (NFL) app to navigate Cowboys Stadium. In 2011, the Tampa Bay Buccaneers became the first NFL club to discontinue the use of paper copies of playbooks, and instead distributed all players their playbook and videos in electronic format via an iPad 2.
The iPad is able to support many music creation applications in addition to the iTunes music playback software. These include sound samplers, guitar and voice effects processors, sequencers for synthesized sounds and sampled loops, virtual synthesizers and drum machines, theremin-style and other touch responsive instruments, drum pads and many more. Gorillaz's 2010 album, The Fall, was created almost exclusively using the iPad by Damon Albarn while on tour with the band. The music video for Luna Sea's 2012 single, "Rouge", was filmed entirely on an iPad.
Due to its popularity, the term "iPad" is occasionally used as a generic name for tablet computers.
Timeline
| Technology | Specific hardware | null |
3910825 | https://en.wikipedia.org/wiki/Cusk-eel | Cusk-eel | The cusk-eel family, Ophidiidae, is a group of marine bony fishes in the Ophidiiformes order. The scientific name is from the Greek ophis meaning "snake", and refers to their eel-like appearance. True eels diverged from other ray-finned fish during the Jurassic, while cusk-eels are part of the Percomorpha clade, along with tuna, perch, seahorses and others.
The oldest fossil cusk-eel is Ampheristus, a highly successful genus with numerous species that existed from the Late Cretaceous (Maastrichtian) to the early Oligocene.
Distribution
Cusk-eels lives in temperate and tropical oceans throughout the world. They live close to the sea bottom, ranging from shallow water to the hadal zone. One species, Abyssobrotula galatheae, was recorded at the bottom of the Puerto Rico Trench, making it the deepest recorded fish at .
Ecology
Cusk-eels are generally very solitary in nature, but some species have been seen to associate themselves with tube worm communities. Liking to be hidden when they are not foraging, they generally associate themselves within muddy bottoms, sinkholes, or larger structures that they can hide in or around, such as caves, coral crevices, or communities of bottom-dwelling invertebrates, with some parasitic species of cusk-eel actually living inside of invertebrate hosts, such as oysters, clams and sea cucumbers. Cusk-eels generally feed nocturnally, preying on invertebrates, crustaceans and other small bottom-dwelling fishes.
Phylogeny
Due to the inconsistencies in specific morphological characteristics in closely related species, attempts to use different characters, such as the position of pelvic fins, to classify Ophiididae into distinct families has proven highly unsatisfactory. Overall, Ophidiidae are classified based on whether or not they practice viviparity and the structures they contain that are associated with bearing life.
Characteristics
Cusk-eels are characterized by a long, slender body that is about 12–13 times as long as it is deep. The largest species, Lamprogrammus shcherbachevi, grows up to in length, but most species are shorter than . Their dorsal and anal fins are typically continuous with the caudal fin (with exception to a few species), forming a long, ribbon like fin around the posterior of the cusk-eel's body. This caudal fin will often be seen to be reduced to a fleshy or bony point, especially when confluent with the dorsal and anal fins. The dorsal fin to anal fin ray ratio is approximately 1.5:1, leading to the dorsal fin typically being longer than the anal. The pectoral fins of cusk-eels are typically longer than the length of their head. Unlike true eels of the order Anguilliformes, cusk-eels have ventral fins that are developed into a forked barbel-like organ below the mouth. In true eels by contrast, the ventral fins are never well-developed and usually missing entirely. Cusk-eels have large mouths relative to their heads, with the upper jaw reaching beyond the eye, and paired nostrils on either side of the head. In cusk-eels, scales are potentially absent; when present, they are small.
Reproduction
Unlike their close relatives, the viviparous brotulas of the family Bythitidae, cusk-eel species are egg-bearing, or oviparous, organisms. While the specifics of the eggs of the family Ophidiidae are unknown, they are believed to be either spawned as individual, free-floating eggs in the open water or are placed in a mucilaginous raft, which will float for several days until they hatch into cusk-eel larvae. These larvae live amongst the plankton relatively close to the water's surface and are believed to control their metamorphoses into adult cusk-eels, dispersing over greater distances into less utilized habitats and reducing competition in concentrated areas.
Conservation status
While a few species are fished commercially – most notably the pink cusk-eel, Genypterus blacodes – and several species of the order Ophidiiformes are listed as vulnerable, not enough information has been gathered about Ophidiidae as a whole to determine their conservation status.
Genera
The cusk-eel family contains about 240 species, grouped into 50 genera:
Genus †Ampheristus (extinct)
Subfamily Brotulinae
Genus Brotula – typical brotulas
Subfamily Brotulotaenilinae
Genus Brotulotaenia
Subfamily Neobythitinae
Genus Abyssobrotula
Genus Acanthonus – boney-eared assfish
Genus Alcockia
Genus Apagesoma
Genus Barathrites
Genus Barathrodemus
Genus Bassogigas
Genus Bassozetus
Genus Bathyonus
Genus Benthocometes
Genus Dannevigia – Australian tusk
Genus Dicrolene
Genus Enchelybrotula
Genus Epetriodus – needletooth cusk
Genus Eretmichthys
Genus Glyptophidium
Genus Holcomycteronus
Genus Homostolus – filament cusk
Genus Hoplobrotula
Genus Hypopleuron – whiptail cusk
Genus Lamprogrammus
Genus Leptobrotula
Genus Leucicorus
Genus Luciobrotula
Genus Mastigopterus
Genus Monomitopus
Genus Neobythites
Genus Neobythitoides
Genus Penopus
Genus Petrotyx
Genus Porogadus
Genus Pycnocraspedum
Genus Selachophidium – Gunther's cusk-eel
Genus Sirembo
Genus Spectrunculus
Genus Spottobrotula
Genus Tauredophidium
Genus Tenuicephalus
Genus Typhlonus
Genus Ventichthys – East-Pacific ventbrotula
Genus Xyelacyba
Subfamily Ophidiinae
Genus Cherublemma – black brotula
Genus Chilara – spotted cusk-eel
Genus Genypterus
Genus Lepophidium
Genus Menziesichthys
Genus Ophidion
Genus Otophidium
Genus Parophidion
Genus Raneya – banded cusk-eel
Gallery
| Biology and health sciences | Acanthomorpha | Animals |
3912709 | https://en.wikipedia.org/wiki/Multiple%20%28mathematics%29 | Multiple (mathematics) | In mathematics, a multiple is the product of any quantity and an integer. In other words, for the quantities a and b, it can be said that b is a multiple of a if b = na for some integer n, which is called the multiplier. If a is not zero, this is equivalent to saying that is an integer.
When a and b are both integers, and b is a multiple of a, then a is called a divisor of b. One says also that a divides b. If a and b are not integers, mathematicians prefer generally to use integer multiple instead of multiple, for clarification. In fact, multiple is used for other kinds of product; for example, a polynomial p is a multiple of another polynomial q if there exists third polynomial r such that p = qr.
Examples
14, 49, −21 and 0 are multiples of 7, whereas 3 and −6 are not. This is because there are integers that 7 may be multiplied by to reach the values of 14, 49, 0 and −21, while there are no such integers for 3 and −6. Each of the products listed below, and in particular, the products for 3 and −6, is the only way that the relevant number can be written as a product of 7 and another real number:
is not an integer;
is not an integer.
Properties
0 is a multiple of every number ().
The product of any integer and any integer is a multiple of . In particular, , which is equal to , is a multiple of (every integer is a multiple of itself), since 1 is an integer.
If and are multiples of then and are also multiples of .
Submultiple
In some texts, "a is a submultiple of b" has the meaning of "a being a unit fraction of b" (ab/n) or, equivalently, "b being an integer multiple n of a" (bna). This terminology is also used with units of measurement (for example by the BIPM and NIST), where a unit submultiple is obtained by prefixing the main unit, defined as the quotient of the main unit by an integer, mostly a power of 103. For example, a millimetre is the 1000-fold submultiple of a metre. As another example, one inch may be considered as a 12-fold submultiple of a foot, or a 36-fold submultiple of a yard.
| Mathematics | Basics | null |
3917234 | https://en.wikipedia.org/wiki/Tetracycline%20antibiotics | Tetracycline antibiotics | Tetracyclines are a group of broad-spectrum antibiotic compounds that have a common basic structure and are either isolated directly from several species of Streptomyces bacteria or produced semi-synthetically from those isolated compounds. Tetracycline molecules comprise a linear fused tetracyclic nucleus (rings designated A, B, C and D) to which a variety of functional groups are attached. Tetracyclines are named after their four ("tetra-") hydrocarbon rings ("-cycl-") derivation ("-ine"). They are defined as a subclass of polyketides, having an octahydrotetracene-2-carboxamide skeleton and are known as derivatives of polycyclic naphthacene carboxamide.
While all tetracyclines have a common structure, they differ from each other by the presence of chloro, methyl, and hydroxyl groups. These modifications do not change their broad antibacterial activity, but do affect pharmacological properties such as half-life and binding to proteins in serum.
Tetracyclines were discovered in the 1940s and exhibited activity against a wide range of microorganisms including gram-positive and gram-negative bacteria, chlamydiota, mycoplasmatota, rickettsiae, and protozoan parasites. Tetracycline itself was discovered later than chlortetracycline and oxytetracycline but is still considered as the parent compound for nomenclature purposes.
Tetracyclines are among the cheapest classes of antibiotics available and have been used extensively in prophylaxis and in treatment of human and animal infections, as well as at subtherapeutic levels in animal feed as growth promoters.
Tetracyclines are growth inhibitors (bacteriostatic) rather than killers of the infectious agent (bacteriocidal) and are only effective against multiplying microorganisms. They are short-acting and passively diffuse through porin channels in the bacterial membrane. They inhibit protein synthesis by binding reversibly to the bacterial 30S ribosomal subunit and preventing the aminoacyl tRNA from binding to the A site of the ribosome. They also bind to some extent the bacterial 50S ribosomal subunit and may alter the cytoplasmic membrane causing intracellular components to leak from bacterial cells.
Tetracyclines all have the same antibacterial spectrum, although there are differences in species' sensitivity to types of tetracyclines. Tetracyclines inhibit protein synthesis in both bacterial and human cells. Bacteria have a system that allows tetracyclines to be transported into the cell, whereas human cells do not. Human cells therefore are spared the effects of tetracycline on protein synthesis.
Tetracyclines retain an important role in medicine, although their usefulness has been reduced with the onset of antibiotic resistance. Tetracyclines remain the treatment of choice for some specific indications.
Because not all of the tetracycline administered orally is absorbed from the gastrointestinal tract, the bacterial population of the intestine can become resistant to tetracyclines, resulting in overgrowth of resistant organisms. The widespread use of tetracyclines is thought to have contributed to an increase in the number of tetracycline-resistant organisms, in turn rendering certain infections more resilient to treatment.
Tetracycline resistance is often due to the acquisition of new genes, which code for energy-dependent efflux of tetracyclines or for a protein that protects bacterial ribosomes from the action of tetracyclines. Furthermore, a limited number of bacteria acquire resistance to tetracyclines by mutations.
Medical uses
Tetracyclines are generally used in the treatment of infections of the urinary tract, respiratory tract, and the intestines and are also used in the treatment of chlamydia, especially in patients allergic to β-lactams and macrolides; however, their use for these indications is less popular than it once was due to widespread development of resistance in the causative organisms.
Tetracyclines are widely used in the treatment of moderately severe acne and rosacea (tetracycline, oxytetracycline, doxycycline or minocycline).
Anaerobic bacteria are not as susceptible to tetracyclines as are aerobic bacteria.
Doxycycline is also used as a prophylactic treatment for infection by Bacillus anthracis (anthrax) and is effective against Yersinia pestis, the infectious agent of bubonic plague. It is also used for malaria treatment and prophylaxis, as well as treating elephantitis filariasis.
Tetracyclines remain the treatment of choice for infections caused by chlamydia (trachoma, psittacosis, salpingitis, urethritis and L. venereum infection), Rickettsia (typhus, Rocky Mountain spotted fever), brucellosis and spirochetal infections (Lyme disease/borreliosis and syphilis). They are also used in veterinary medicine.
They may have a role in reducing the duration and severity of cholera, although drug-resistance is mounting and their effect on overall mortality is questioned.
Side effects
Side-effects from tetracyclines are not common, but of particular note is phototoxicity. It increases the risk of sunburn under exposure to light from the sun or other sources. This may be of particular importance for those intending to take on vacations long-term doxycycline as a malaria prophylaxis.
They may cause stomach or bowel upsets, and, on rare occasions, allergic reactions. Very rarely, severe headache and vision problems may be signs of dangerous secondary intracranial hypertension, also known as idiopathic intracranial hypertension.
Tetracyclines are teratogens due to the likelihood of causing teeth discolouration in the fetus as they develop in infancy. For this same reason, tetracyclines are contraindicated for use in children under 8 years of age. Some adults also experience teeth discoloration (mild grey hue) after use. They are, however, safe to use in the first 18 weeks of pregnancy.
Some patients taking tetracyclines require medical supervision because they can cause steatosis and liver toxicity.
Cautions
Tetracyclines should be used with caution by those with liver impairment. Also, because the molecules are soluble in water it can worsen kidney failure (this is not true of the lipid-soluble agents doxycycline and minocycline). They may increase muscle weakness in myasthenia gravis and exacerbate systemic lupus erythematosus. Antacids containing aluminium and calcium reduce the absorption of all tetracyclines, and dairy products reduce absorption greatly for all but minocycline.
The breakdown products of tetracyclines are toxic and can cause Fanconi syndrome, a potentially fatal disease affecting proximal tubular function in the nephrons of the kidney. Prescriptions of these drugs should be discarded once expired because they can cause hepatotoxicity.
It was once believed that tetracycline antibiotics impair the effectiveness of many types of hormonal contraception. Recent research has shown no significant loss of effectiveness in oral contraceptives while using most tetracyclines. Despite these studies, many physicians still recommend the use of barrier contraception for people taking any tetracyclines to prevent unwanted pregnancy.
In tetracycline preparation, stability must be considered in order to avoid formation of toxic epi-anhydrotetracyclines.
Contraindications
Tetracycline use should be avoided in pregnant or lactating women, and in children with developing teeth because they may result in permanent staining (dark yellow-gray teeth with a darker horizontal band that goes across the top and bottom rows of teeth), and possibly affect the growth of teeth and bones.
Usage during the first 12 weeks of pregnancy does not appear to increase the risk of any major birth defects. There may be a small increased risk for minor birth defects such as an inguinal hernia, but the number of reports is too small to be sure if there actually is any risk.
Mechanism of action
Tetracycline antibiotics are protein synthesis inhibitors. They inhibit the initiation of translation in variety of ways by binding to the 30S ribosomal subunit, which is made up of 16S rRNA and 21 proteins. They inhibit the binding of aminoacyl-tRNA to the mRNA translation complex. Some studies have shown that tetracyclines may bind to both 16S and 23S rRNAs. Tetracyclines also have been found to inhibit matrix metalloproteinases. This mechanism does not add to their antibiotic effects, but has led to extensive research on chemically modified tetracyclines or CMTs (like incyclinide) for the treatment of rosacea, acne, diabetes and various types of neoplasms.
It has been shown that tetracyclines are not only active against broad spectrum of bacteria, but also against viruses, protozoa that lack mitochondria and some noninfectious conditions. The binding of tetracyclines to cellular dsRNA (double stranded RNA) may be an explanation for their wide range of effect. It can also be attributed to the nature of ribosomal protein synthesis pathways among bacteria.
Incyclinide was announced to be ineffective for rosacea in September 2007.
Several trials have examined modified and unmodified tetracyclines for the treatment of human cancers; of those, very promising results were achieved with CMT-3 for patients with Kaposi Sarcoma.
Structure-activity relationship
Tetracyclines are composed of a rigid skeleton of 4 fused rings. The rings structure of tetracyclines is divided into an upper modifiable region and a lower non modifiable region. An active tetracycline requires a C10 phenol as well as a C11-C12 keto-enol substructure in conjugation with a 12a-OH group and a C1-C3 diketo substructure. Removal of the dimethylamine group at C4 reduces antibacterial activity. Replacement of the carboxylamine group at C2 results in reduced antibacterial activity but it is possible to add substituents to the amide nitrogen to get more soluble analogs like the prodrug lymecycline. The simplest tetracycline with measurable antibacterial activity is 6-deoxy-6-demethyltetracycline and its structure is often considered to be the minimum pharmacophore for the tetracycle class of antibiotics. C5-C9 can be modified to make derivatives with varying antibacterial activity.
Mechanism of resistance
Cells can become resistant to tetracycline by enzymatic inactivation of tetracycline, efflux, ribosomal protection, reduced permeability and ribosome mutation.
Inactivation is the rarest type of resistance, where NADPH-dependent oxidoreductase, a class of antibiotic destructase, modifies the tetracycline antibiotic at their oxidative soft spot leading to an inactivation of the tetracycline antibiotic. For example, the oxireductase makes a modification on the C11a site of oxytetracycline. Both Mg2+ chelation and ribosome binding are required for the biological activity of oxytetracycline and the modification attenuate the binding, leading to inactivation of the oxytetracycline antibiotic.
In the most common mechanism of reaction, efflux, various resistance genes encode a membrane protein that actively pumps tetracycline out of the cell by exchanging a proton for a tetracycline cation complex. This exchange leads to a reduced cytoplasmic concentration of tetracycline.
In ribosomal protection, a resistance gene encodes a protein that can have several effects, depending on what gene is transferred. Twelve classes of ribosomal protection genes/proteins have been found.
Possible mechanisms of action of these protective proteins include:
blocking tetracyclines from binding to the ribosome
binding to the ribosome and distorting the structure to still allow t-RNA binding while tetracycline is bound
binding to the ribosome and dislodging tetracycline
Administration
When ingested, it is usually recommended that the more water-soluble, short-acting tetracyclines (plain tetracycline, chlortetracycline, oxytetracycline, demeclocycline and methacycline) be taken with a full glass of water, either two hours after eating or two hours before eating. This is partly because most tetracyclines bind with food and also easily with magnesium, aluminium, iron and calcium, which reduces their ability to be completely absorbed by the body. Dairy products, antacids and preparations containing iron should be avoided near the time of taking the drug. Partial exceptions to these rules occur for doxycycline and minocycline, which may be taken with food (though not iron, antacids, or calcium supplements). Minocycline can be taken with dairy products because it does not chelate calcium as readily, although dairy products do decrease absorption of minocycline slightly.
History
The history of the tetracyclines involves the collective contributions of thousands of dedicated researchers, scientists, clinicians, and business executives. Tetracyclines were discovered in the 1940s, first reported in scientific literature in 1948, and exhibited activity against a wide range of microorganisms. The first members of the tetracycline group to be described were chlortetracycline and oxytetracycline. Chlortetracycline (Aureomycin) was first discovered as an ordinary item in 1945 and initially endorsed in 1948 by Benjamin Minge Duggar, a 73-year-old emeritus professor of botany employed by American Cyanamid – Lederle Laboratories, under the leadership of Yellapragada Subbarow. Duggar derived the substance from a Missouri soil sample, golden-colored, fungus-like, soil-dwelling bacterium named Streptomyces aureofaciens.
About the same time as Lederle discovered aureomycin, Pfizer was scouring the globe for new antibiotics. Soil samples were collected from jungles, deserts, mountaintops, and oceans. But ultimately oxytetracycline (terramycin) was isolated in 1949 by Alexander Finlay from a soil sample collected on the grounds of a factory in Terre Haute, Indiana. It came from a similar soil bacterium named Streptomyces rimosus. From the beginning, terramycin was a molecule enveloped in controversy. It was the subject of the first mass-marketing campaign by a modern pharmaceutical company. Pfizer advertised the drug heavily in medical journals, eventually spending twice as much on marketing as it did to discover and develop terramycin. Still, it turned Pfizer, then a small company, into a pharmaceutical giant.
The Pfizer group, led by Francis A. Hochstein, in loose collaboration with and Robert Burns Woodward, determined the structure of oxytetracycline, enabling Lloyd H. Conover to successfully produce tetracycline itself as a synthetic product. In 1955, Conover discovered that hydrogenolysis of aureomycin gives a deschloro product that is just as active as the original product. This proved for the first time that chemically modified antibiotics could have biological activity. Within a few years, a number of semisynthetic tetracyclines had entered the market, and now most antibiotic discoveries are of novel active derivatives of older compounds.
Other tetracyclines were identified later, either as naturally occurring molecules, e.g., tetracycline from S. aureofaciens, S. rimosus, and S. viridofaciens and dimethyl-chlortetracycline from S. aureofaciens, or as products of semisynthetic approaches, e.g., methacycline, doxycycline, and minocycline.
Research conducted by anthropologist George J. Armelagos and his team at Emory University showed that ancient Nubians from the post-Meroitic period (around AD 350) had deposits of tetracycline in their bones, detectable through analyses of cross-sections through ultraviolet light – the deposits are fluorescent, just as are modern ones. Armelagos suggested that this was due to ingestion of the local ancient beer (very much like the Egyptian beer), made from contaminated stored grains.
Development
Tetracyclines were noted for their broad spectrum antibacterial activity and were commercialized with clinical success beginning in the late 1940s to the early 1950s. The second-generation semisynthetic analogs and more recent third-generation compounds show the continued evolution of the tetracycline platform towards derivatives with increased potency as well as efficacy against tetracycline-resistant bacteria, with improved pharmacokinetic and chemical properties.
Shortly after the introduction of tetracycline therapy, the first tetracycline-resistant bacterial pathogen was identified. Since then, tetracycline-resistant bacterial pathogens have continued to be identified, limiting tetracycline's effectiveness in treatment of bacterial disease.
Glycylcyclines and fluorocyclines are new classes of antibiotics derived from tetracycline. These tetracycline analogues are specifically designed to overcome two common mechanisms of tetracycline resistance, namely resistance mediated by acquired efflux pumps and/or ribosomal protection.
In 2005, tigecycline, the first member of a new subgroup of tetracyclines named glycylcyclines, was introduced to treat infections that are resistant to other antimicrobials. Although it is structurally related to minocycline, alterations to the molecule resulted in its expanded spectrum of activity and decreased susceptibility to the development of resistance when compared with other tetracycline antibiotics. Like minocycline, tigecycline binds to the bacterial 30S ribosome, blocking the entry of transfer RNA. This ultimately prevents protein synthesis and thus inhibiting bacterial growth. However, the addition of an N,N,-dimethylglycylamido group at the 9 position of the minocycline molecule increases the affinity of tigecycline for the ribosomal target up to 5 times when compared with minocycline or tetracycline. This allows for an expanded spectrum of activity and decreased susceptibility to the development of resistance. While tigecycline was the first tetracycline approved in over 20 years, other, newer versions of tetracyclines are currently in human clinical trials.
List of tetracycline antibiotics
Use as research reagents
Members of the tetracycline class of antibiotics are often used as research reagents in in vitro and in vivo biomedical research experiments involving bacteria as well in experiments in eukaryotic cells and organisms with inducible protein expression systems using tetracycline-controlled transcriptional activation. The mechanism of action for the antibacterial effect of tetracyclines relies on disrupting protein translation in bacteria, thereby damaging the ability of microbes to grow and repair; however protein translation is also disrupted in eukaryotic mitochondria leading to effects that may confound experimental results.
It can be used as an artificial biomarker in wildlife to check if wild animals are consuming a bait that contains a vaccine or medication. Since it is fluorescent and binds to calcium, a UV lamp can be used to check if it is in a tooth pulled from an animal. For example, it was used to check uptake of oral rabies vaccine baits by raccoons in the USA. However, this is an invasive procedure for the animal and labour-intensive for the researcher. Therefore, other dyes such as rhodamine B that can be detected in hair and whiskers are preferred.
| Biology and health sciences | Antibiotics | Health |
3917813 | https://en.wikipedia.org/wiki/Rutherford%20model | Rutherford model | The Rutherford model is a name for the first model of an atom with a compact nucleus. The concept arose from Ernest Rutherford discovery of the nucleus. Rutherford directed the Geiger–Marsden experiment in 1909, which showed much more alpha particle recoil than J. J. Thomson's plum pudding model of the atom could explain. Thomson's model had positive charge spread out in the atom. Rutherford's analysis proposed a high central charge concentrated into a very small volume in comparison to the rest of the atom and with this central volume containing most of the atom's mass. The central region would later be known as the atomic nucleus. Rutherford did not discuss the organization of electrons in the atom and did not himself propose a model for the atom. Niels Bohr joined Rutherford's lab and developed a theory for the electron motion which became known as the Bohr model.
Background
Throughout the 1800's speculative ideas about atoms were discussed and published.
JJ Thomson's model was the first of these models to be based on experimentally detected subatomic particles. In the same paper that Thomson announced his results on "corpuscle" nature of cathode rays, an event considered the discovery of the electron, he began speculating on atomic models composed of electrons. He developed his model, now called the plum pudding model, primarily in 1904-06. He produced an elaborate mechanical model of the electrons moving in concentric rings, but the positive charge needed to balance the negative electrons was a simple sphere of uniform charge and unknown composition. Between 1904 and 1910 Thomson developed formulae for the deflection of fast beta particles from his atomic model for comparison to experiment. Similar work by Rutherford using alpha particles would eventually show Thomson's model could not be correct.
Also among the early models were "planetary" or Solar System-like models. In a 1901 paper, Jean Baptiste Perrin used Thomson's discovery in a proposed a Solar System like model for atoms, with very strongly charged "positive suns" surrounded by "corpuscles, a kind of small negative planets", where the word "corpuscles" refers to what we now call electrons. Perrin discussed how this hypothesis might related to important then unexplained phenomena like the photoelectric effect, emission spectra, and radioactivity. Perrin later credited Rutherford with the discovery of the nuclear model.
A somewhat similar model proposed by Hantaro Nagaoka in 1904 used Saturn's rings as an analog. The rings consisted of a large number of particles that repelled each other but were attracted to a large central charge. This charge was calculated to be 10,000 times the charge of the ring particles for stability. George A. Schott showed in 1904 that Nagaoka's model could not be consistent with results of atomic spectroscopy and the model fell out of favor.
Experimental basis for the model
Rutherford's nuclear model of the atom grew out of a series of experiments with alpha particles, a form of radiation Rutherford discovered in 1899. These experiments demonstrated that alpha particles "scattered" or bounced off atoms in ways unlike Thomson's model predicted. In 1908 and 1910, Hans Geiger and Ernest Marsden in Rutherford's lab showed that alpha particles could occasionally be reflected from gold foils. If Thomson was correct, the beam would go through the gold foil with very small deflections. In the experiment most of the beam passed through the foil, but a few were deflected.
In a May 1911 paper, Rutherford presented his own physical model for subatomic structure, as an interpretation for the unexpected experimental results. In it, the atom is made up of a central charge (this is the modern atomic nucleus, though Rutherford did not use the term "nucleus" in his paper). Rutherford only committed himself to a small central region of very high positive or negative charge in the atom.
For concreteness, consider the passage of a high speed α particle through an atom having a positive central charge N e, and surrounded by a compensating charge of N electrons.
Using only energetic considerations of how far particles of known speed would be able to penetrate toward a central charge of 100 e, Rutherford was able to calculate that the radius of his gold central charge would need to be less (how much less could not be told) than 3.4 × 10−14 meters. This was in a gold atom known to be 10−10 metres or so in radius—a very surprising finding, as it implied a strong central charge less than 1/3000th of the diameter of the atom.
The Rutherford model served to concentrate a great deal of the atom's charge and mass to a very small core, but did not attribute any structure to the remaining electrons and remaining atomic mass. It did mention the atomic model of Hantaro Nagaoka, in which the electrons are arranged in one or more rings, with the specific metaphorical structure of the stable rings of Saturn. The plum pudding model of J. J. Thomson also had rings of orbiting electrons.
The Rutherford paper suggested that the central charge of an atom might be "proportional" to its atomic mass in hydrogen mass units u (roughly 1/2 of it, in Rutherford's model). For gold, this mass number is 197 (not then known to great accuracy) and was therefore modelled by Rutherford to be possibly 196 u. However, Rutherford did not attempt to make the direct connection of central charge to atomic number, since gold's "atomic number" (at that time merely its place number in the periodic table) was 79, and Rutherford had modelled the charge to be about +100 units (he had actually suggested 98 units of positive charge, to make half of 196). Thus, Rutherford did not formally suggest the two numbers (periodic table place, 79, and nuclear charge, 98 or 100) might be exactly the same.
In 1913 Antonius van den Broek suggested that the nuclear charge and atomic weight were not connected, clearing the way for the idea that atomic number and nuclear charge were the same. This idea was quickly taken up by Rutherford's team and was confirmed experimentally within two years by Henry Moseley.
These are the key indicators:
The atom's electron cloud does not (substantially) influence alpha particle scattering.
Much of an atom's positive charge is concentrated in a relatively tiny volume at the center of the atom, known today as the nucleus. The magnitude of this charge is proportional to (up to a charge number that can be approximately half of) the atom's atomic mass—the remaining mass is now known to be mostly attributed to neutrons. This concentrated central mass and charge is responsible for deflecting both alpha and beta particles.
The mass of heavy atoms such as gold is mostly concentrated in the central charge region, since calculations show it is not deflected or moved by the high speed alpha particles, which have very high momentum in comparison to electrons, but not with regard to a heavy atom as a whole.
The atom itself is about 100,000 (105) times the diameter of the nucleus. This could be related to putting a grain of sand in the middle of a football field.
Contribution to modern science
Rutherford's new atom model caused no reaction at first. Rutherford explicitly ignores the electrons, only mentioning Hantaro Nagaoka's Saturnian model. By ignoring the electrons Rutherford also ignores any potential implications for atomic spectroscopy for chemistry. Rutherford himself did not press the case for his atomic model in the following years: his own 1913 book on "Radioactive substances and their radiations" only mentions the atom twice; other books by other authors around this time focus on Thomson's model.
The impact of Rutherford's nuclear model came after Niels Bohr arrived as a post-doctoral student in Manchester at Rutherford's invitation. Bohr dropped his work on the Thomson model in favor of Rutherford's nuclear model, developing the Rutherford–Bohr model over the next several years. Eventually Bohr incorporated early ideas of quantum mechanics into the model of the atom, allowing prediction of electronic spectra and concepts of chemistry.
After Rutherford's discovery, subsequent research determined the atomic structure which led to Rutherford's gold foil experiment. Scientists eventually discovered that atoms have a positively charged nucleus (with an atomic number of charges) in the center, with a radius of about 1.2 × 10−15 meters × [atomic mass number]. Electrons were found to be even smaller.
| Physical sciences | Atomic physics | Physics |
21578485 | https://en.wikipedia.org/wiki/Color%20space | Color space | A color space is a specific organization of colors. In combination with color profiling supported by various physical devices, it supports reproducible representations of colorwhether such representation entails an analog or a digital representation. A color space may be arbitrary, i.e. with physically realized colors assigned to a set of physical color swatches with corresponding assigned color names (including discrete numbers infor examplethe Pantone collection), or structured with mathematical rigor (as with the NCS System, Adobe RGB and sRGB). A "color space" is a useful conceptual tool for understanding the color capabilities of a particular device or digital file. When trying to reproduce color on another device, color spaces can show whether shadow/highlight detail and color saturation can be retained, and by how much either will be compromised.
A "color model" is an abstract mathematical model describing the way colors can be represented as tuples of numbers (e.g. triples in RGB or quadruples in CMYK); however, a color model with no associated mapping function to an absolute color space is a more or less arbitrary color system with no connection to any globally understood system of color interpretation. Adding a specific mapping function between a color model and a reference color space establishes within the reference color space a definite "footprint", known as a gamut, and for a given color model, this defines a color space. For example, Adobe RGB and sRGB are two different absolute color spaces, both based on the RGB color model. When defining a color space, the usual reference standard is the CIELAB or CIEXYZ color spaces, which were specifically designed to encompass all colors the average human can see.
Since "color space" identifies a particular combination of the color model and the mapping function, the word is often used informally to identify a color model. However, even though identifying a color space automatically identifies the associated color model, this usage is incorrect in a strict sense. For example, although several specific color spaces are based on the RGB color model, there is no such thing as the singular RGB color space.
History
In 1802, Thomas Young postulated the existence of three types of photoreceptors (now known as cone cells) in the eye, each of which was sensitive to a particular range of visible light. Hermann von Helmholtz developed the Young–Helmholtz theory further in 1850: that the three types of cone photoreceptors could be classified as short-preferring (blue), middle-preferring (green), and long-preferring (red), according to their response to the wavelengths of light striking the retina. The relative strengths of the signals detected by the three types of cones are interpreted by the brain as a visible color. But it is not clear that they thought of colors as being points in color space.
The color-space concept was likely due to Hermann Grassmann, who developed it in two stages. First, he developed the idea of vector space, which allowed the algebraic representation of geometric concepts in n-dimensional space. Fearnley-Sander (1979) describes Grassmann's foundation of linear algebra as follows:
With this conceptual background, in 1853, Grassmann published a theory of how colors mix; it and its three color laws are still taught, as Grassmann's law.
Examples
Colors can be created in printing with color spaces based on the CMYK color model, using the subtractive primary colors of pigment (cyan, magenta, yellow, and black). To create a three-dimensional representation of a given color space, we can assign the amount of magenta color to the representation's X axis, the amount of cyan to its Y axis, and the amount of yellow to its Z axis. The resulting 3-D space provides a unique position for every possible color that can be created by combining those three pigments.
Colors can be created on computer monitors with color spaces based on the RGB color model, using the additive primary colors (red, green, and blue). A three-dimensional representation would assign each of the three colors to the X, Y, and Z axes. Colors generated on a given monitor will be limited by the reproduction medium, such as the phosphor (in a CRT monitor) or filters and backlight (LCD monitor).
Another way of creating colors on a monitor is with an HSL or HSV color model, based on hue, saturation, brightness (value/lightness). With such a model, the variables are assigned to cylindrical coordinates.
Many color spaces can be represented as three-dimensional values in this manner, but some have more, or fewer dimensions, and some, such as Pantone, cannot be represented in this way at all.
Conversion
Color space conversion is the translation of the representation of a color from one basis to another. This typically occurs in the context of converting an image that is represented in one color space to another color space, the goal being to make the translated image look as similar as possible to the original.
RGB density
The RGB color model is implemented in different ways, depending on the capabilities of the system used. The most common incarnation in general use is the 24-bit implementation, with 8 bits, or 256 discrete levels of color per channel. Any color space based on such a 24-bit RGB model is thus limited to a range of 256×256×256 ≈ 16.7 million colors. Some implementations use 16 bits per component for 48 bits total, resulting in the same gamut with a larger number of distinct colors. This is especially important when working with wide-gamut color spaces (where most of the more common colors are located relatively close together), or when a large number of digital filtering algorithms are used consecutively. The same principle applies for any color space based on the same color model, but implemented at different bit depths.
Lists
CIE 1931 XYZ color space was one of the first attempts to produce a color space based on measurements of human color perception (earlier efforts were by James Clerk Maxwell, König & Dieterici, and Abney at Imperial College) and it is the basis for almost all other color spaces. The CIERGB color space is a linearly-related companion of CIE XYZ. Additional derivatives of CIE XYZ include the CIELUV, CIEUVW, and CIELAB.
Generic
RGB uses additive color mixing, because it describes what kind of light needs to be emitted to produce a given color. RGB stores individual values for red, green and blue. RGBA is RGB with an additional channel, alpha, to indicate transparency. Common color spaces based on the RGB model include sRGB, Adobe RGB, ProPhoto RGB, scRGB, and CIE RGB.
CMYK uses subtractive color mixing used in the printing process, because it describes what kind of inks need to be applied so the light reflected from the substrate and through the inks produces a given color. One starts with a white substrate (canvas, page, etc.), and uses ink to subtract color from white to create an image. CMYK stores ink values for cyan, magenta, yellow and black. There are many CMYK color spaces for different sets of inks, substrates, and press characteristics (which change the dot gain or transfer function for each ink and thus change the appearance).
YIQ was formerly used in NTSC (North America, Japan and elsewhere) television broadcasts for historical reasons. This system stores a luma value roughly analogous to (and sometimes incorrectly identified as) luminance, along with two chroma values as approximate representations of the relative amounts of blue and red in the color. It is similar to the YUV scheme used in most video capture systems and in PAL (Australia, Europe, except France, which uses SECAM) television, except that the YIQ color space is rotated 33° with respect to the YUV color space and the color axes are swapped. The YDbDr scheme used by SECAM television is rotated in another way.
YPbPr is a scaled version of YUV. It is most commonly seen in its digital form, YCbCr, used widely in video and image compression schemes such as MPEG and JPEG.
xvYCC is a new international digital video color space standard published by the IEC (IEC 61966-2-4). It is based on the ITU BT.601 and BT.709 standards but extends the gamut beyond the R/G/B primaries specified in those standards.
HSV (hue, saturation, value), also known as HSB (hue, saturation, brightness) is often used by artists because it is often more natural to think about a color in terms of hue and saturation than in terms of additive or subtractive color components. HSV is a transformation of an RGB color space, and its components and colorimetry are relative to the RGB color space from which it was derived.
HSL (hue, saturation, lightness/luminance), also known as HLS or HSI (hue, saturation, intensity) is quite similar to HSV, with "lightness" replacing "brightness". The difference is that the brightness of a pure color is equal to the brightness of white, while the lightness of a pure color is equal to the lightness of a medium gray.
Commercial
Munsell color system
Pantone Matching System (PMS)
Natural Color System (NCS)
Special-purpose
The RG Chromaticity space is used in computer vision applications. It shows the color of light (red, yellow, green etc.), but not its intensity (dark, bright).
The TSL color space (Tint, Saturation and Luminance) is used in face detection.
Obsolete
Early color spaces had two components. They largely ignored blue light because the added complexity of a 3-component process provided only a marginal increase in fidelity when compared to the jump from monochrome to 2-component color.
RG for early Technicolor film
RGK for early color printing
Absolute color space
In color science, there are two meanings of the term absolute color space:
A color space in which the perceptual difference between colors is directly related to distances between colors as represented by points in the color space, i.e. a uniform color space.
A color space in which colors are unambiguous, that is, where the interpretations of colors in the space are colorimetrically defined without reference to external factors.
In this article, we concentrate on the second definition.
CIEXYZ, sRGB, and ICtCp are examples of absolute color spaces, as opposed to a generic RGB color space.
A non-absolute color space can be made absolute by defining its relationship to absolute colorimetric quantities. For instance, if the red, green, and blue colors in a monitor are measured exactly, together with other properties of the monitor, then RGB values on that monitor can be considered as absolute. The CIE 1976 L*, a*, b* color space is sometimes referred to as absolute, though it also needs a white point specification to make it so.
A popular way to make a color space like RGB into an absolute color is to define an ICC profile, which contains the attributes of the RGB. This is not the only way to express an absolute color, but it is the standard in many industries. RGB colors defined by widely accepted profiles include sRGB and Adobe RGB. The process of adding an ICC profile to a graphic or document is sometimes called tagging or embedding; tagging, therefore, marks the absolute meaning of colors in that graphic or document.
Conversion errors
A color in one absolute color space can be converted into another absolute color space, and back again, in general; however, some color spaces may have gamut limitations, and converting colors that lie outside that gamut will not produce correct results. There are also likely to be rounding errors, especially if the popular range of only 256 distinct values per component (8-bit color) is used.
One part of the definition of an absolute color space is the viewing conditions. The same color, viewed under different natural or artificial lighting conditions, will look different. Those involved professionally with color matching may use viewing rooms, lit by standardized lighting.
Occasionally, there are precise rules for converting between non-absolute color spaces. For example, HSL and HSV spaces are defined as mappings of RGB. Both are non-absolute, but the conversion between them should maintain the same color. However, in general, converting between two non-absolute color spaces (for example, RGB to CMYK) or between absolute and non-absolute color spaces (for example, RGB to L*a*b*) is almost a meaningless concept.
Arbitrary spaces
A different method of defining absolute color spaces is familiar to many consumers as the swatch card, used to select paint, fabrics, and the like. This is a way of agreeing a color between two parties. A more standardized method of defining absolute colors is the Pantone Matching System, a proprietary system that includes swatch cards and recipes that commercial printers can use to make inks that are a particular color.
| Physical sciences | Basics | Physics |
35869003 | https://en.wikipedia.org/wiki/Human%20mouth | Human mouth | In human anatomy, the mouth is the first portion of the alimentary canal that receives food and produces saliva. The oral mucosa is the mucous membrane epithelium lining the inside of the mouth.
In addition to its primary role as the beginning of the digestive system, the mouth also plays a significant role in communication. While primary aspects of the voice are produced in the throat, the tongue, lips, and jaw are also needed to produce the range of sounds included in speech.
The mouth consists of two regions, the vestibule and the oral cavity proper. The mouth, normally moist, is lined with a mucous membrane, and contains the teeth. The lips mark the transition from mucous membrane to skin, which covers most of the body.
Structure
Oral cavity
The mouth consists of two regions: the vestibule and the oral cavity proper. The vestibule is the area between the teeth, lips and cheeks. The oral cavity is bounded at the sides and in front by the alveolar process (containing the teeth) and at the back by the isthmus of the fauces. Its roof is formed by the hard palate. The floor is formed by the mylohyoid muscles and is occupied mainly by the anterior two-thirds of the tongue. A mucous membrane – the oral mucosa, lines the sides and under surface of the tongue to the gums, and lines the inner aspect of the jaw (mandible). It receives secretions from the submandibular and sublingual salivary glands. The posterior border of the oral cavity (ie, junction between the oral cavity and the oropharynx) includes the junction of the hard palate and the soft palate superiorly, the circumvallate papillae of the tongue inferiorly, and the retromolar trigone.
Lips
The lips come together to close the opening of the mouth, forming a line between the upper and lower lip. In facial expression, this mouth line is iconically shaped like an up-open parabola in a smile, and like a down-open parabola in a frown. A down-turned mouth means a mouth line forming a down-turned parabola, and when permanent can be normal. Also, a down-turned mouth can be part of the presentation of Prader–Willi syndrome.
Nerve supply
The teeth and the periodontium (the tissues that support the teeth) are innervated by the maxillary and mandibular nerves – divisions of the trigeminal nerve. Maxillary (upper) teeth and their associated periodontal ligament are innervated by the superior alveolar nerves, branches of the maxillary division, termed the posterior superior alveolar nerve, anterior superior alveolar nerve, and the variably present middle superior alveolar nerve. These nerves form the superior dental plexus above the maxillary teeth. The mandibular (lower) teeth and their associated periodontal ligament are innervated by the inferior alveolar nerve, a branch of the mandibular division. This nerve runs inside the mandible, within the inferior alveolar canal below the mandibular teeth, giving off branches to all the lower teeth (inferior dental plexus). The oral mucosa of the gingiva (gums) on the facial (labial) aspect of the maxillary incisors, canines and premolar teeth is innervated by the superior labial branches of the infraorbital nerve. The posterior superior alveolar nerve supplies the gingiva on the facial aspect of the maxillary molar teeth. The gingiva on the palatal aspect of the maxillary teeth is innervated by the greater palatine nerve apart from in the incisor region, where it is the nasopalatine nerve (long sphenopalatine nerve). The gingiva of the lingual aspect of the mandibular teeth is innervated by the sublingual nerve, a branch of the lingual nerve. The gingiva on the facial aspect of the mandibular incisors and canines is innervated by the mental nerve, the continuation of the inferior alveolar nerve emerging from the mental foramen. The gingiva of the buccal (cheek) aspect of the mandibular molar teeth is innervated by the buccal nerve (long buccal nerve).
Development
The philtrum is the vertical depression formed between the philtral ridges between the upper lip and the nasal septum, formed where the nasomedial and maxillary processes meet during embryo development. When these processes fail to fuse fully, a cleft lip, cleft palate, or both can result.
The nasolabial folds are the deep creases of tissue that extend from the nose to the sides of the mouth. One of the first signs of age on the human face is the increase in prominence of the nasolabial folds.
Function
The mouth plays an important role in eating, drinking, and speaking. Mouth breathing refers to the act of breathing through the mouth (as a temporary backup system) if there is an obstruction to breathing through the nose, which is the designated breathing organ for the human body.
Infants are born with a sucking reflex, by which they instinctively know to suck for nourishment using their lips and jaw.
The mouth also helps in chewing and biting food.
For some disabled people, especially many disabled artists, who through illness, accident or congenital disability have lost dexterity, their mouths take the place of their hands, when typing, texting, writing, making drawings, paintings and other works of art by maneuvering brushes and other tools, in addition to the basic oral functions. Mouth painters hold the brush in their mouth or between their teeth and maneuver it with their tongue and cheek muscles, but mouth painting can be strenuous for neck and jaw muscles since the head has to perform the same back and forth movement as a hand does when painting.
A male mouth can hold, on average, , while a female mouth holds .
| Biology and health sciences | Human anatomy | Health |
2916607 | https://en.wikipedia.org/wiki/Force%20field%20%28physics%29 | Force field (physics) | In physics, a force field is a vector field corresponding with a non-contact force acting on a particle at various positions in space. Specifically, a force field is a vector field , where is the force that a particle would feel if it were at the position .
Examples
Gravity is the force of attraction between two objects. A gravitational force field models this influence that a massive body (or more generally, any quantity of energy) extends into the space around itself. In Newtonian gravity, a particle of mass M creates a gravitational field , where the radial unit vector points away from the particle. The gravitational force experienced by a particle of light mass m, close to the surface of Earth is given by , where g is Earth's gravity.
An electric field exerts a force on a point charge q, given by .
In a magnetic field , a point charge moving through it experiences a force perpendicular to its own velocity and to the direction of the field, following the relation: .
Work
Work is dependent on the displacement as well as the force acting on an object. As a particle moves through a force field along a path C, the work done by the force is a line integral:
This value is independent of the velocity/momentum that the particle travels along the path.
Conservative force field
For a conservative force field, it is also independent of the path itself, depending only on the starting and ending points. Therefore, the work for an object travelling in a closed path is zero, since its starting and ending points are the same:
If the field is conservative, the work done can be more easily evaluated by realizing that a conservative vector field can be written as the gradient of some scalar potential function:
The work done is then simply the difference in the value of this potential in the starting and end points of the path. If these points are given by x = a and x = b, respectively:
| Physical sciences | Classical mechanics | Physics |
2916631 | https://en.wikipedia.org/wiki/Sodium%20sulfide | Sodium sulfide | Sodium sulfide is a chemical compound with the formula Na2S, or more commonly its hydrate Na2S·9H2O. Both the anhydrous and the hydrated salts in pure crystalline form are colorless solids, although technical grades of sodium sulfide are generally yellow to brick red owing to the presence of polysulfides and commonly supplied as a crystalline mass, in flake form, or as a fused solid. They are water-soluble, giving strongly alkaline solutions. When exposed to moisture, Na2S immediately hydrates to give sodium hydrosulfide.
Some commercial samples are specified as Na2S·xH2O, where a weight percentage of Na2S is specified. Commonly available grades have around 60% Na2S by weight, which means that x is around 3. These grades of sodium sulfide are often marketed as 'sodium sulfide flakes'. These samples consist of NaSH, NaOH, and water.
Structure
The structures of sodium sulfides have been determined by X-ray crystallography. The nonahydrate features S2- hydrogen-bonded to 12 water molecules. The pentahydrate consists of S2- centers bound to Na+ and encased by an array of hydrogen bonds. Anhydrous Na2S, which is rarely encountered, adopts the antifluorite structure, which means that the Na+ centers occupy sites of the fluoride in the CaF2 framework, and the larger S2− occupy the sites for Ca2+.
Production
Industrially Na2S is produced by carbothermic reduction of sodium sulfate often using coal:
Na2SO4 + 2 C → Na2S + 2 CO2
In the laboratory, the salt can be prepared by reduction of sulfur with sodium in anhydrous ammonia, or by sodium in dry THF with a catalytic amount of naphthalene (forming sodium naphthalenide):
2 Na + S → Na2S
Reactions with inorganic reagents
The sulfide ion in sulfide salts such as sodium sulfide can incorporate a proton into the salt by protonation:
+ →
Because of this capture of the proton (), sodium sulfide has basic character. Sodium sulfide is strongly basic, able to absorb two protons. Its conjugate acid is sodium hydrosulfide (). An aqueous solution contains a significant portion of sulfide ions that are singly protonated.
+ +
+ +
Sodium sulfide is unstable in the presence of water due to the gradual loss of hydrogen sulfide into the atmosphere.
When heated with oxygen and carbon dioxide, sodium sulfide can oxidize to sodium carbonate and sulfur dioxide:
2 Na2S + 3 O2 + 2 → 2 Na2CO3 + 2 SO2
Oxidation with hydrogen peroxide gives sodium sulfate:
Na2S + 4 H2O2 → 4 + Na2SO4
Upon treatment with sulfur, sodium polysulfides are formed:
2 Na2S + S8 → 2 Na2S5
Pulp and paper industry
In terms of its dominant use, "sodium sulfide" is primarily used in the kraft process in the pulp and paper industry. It aids in the delignification process, affording cellulose, which is the main component of paper.
It is used in water treatment as an oxygen scavenger agent and also as a metals precipitant; in chemical photography for toning black and white photographs; in the textile industry as a bleaching agent, for desulfurising and as a dechlorinating agent; and in the leather trade for the sulfitisation of tanning extracts. It is used in chemical manufacturing as a sulfonation and sulfomethylation agent. It is used in the production of rubber chemicals, sulfur dyes and other chemical compounds. It is used in other applications including ore flotation, oil recovery, making dyes, and detergent. It is also used during leather processing, as an unhairing agent in the liming operation.
Reagent in organic chemistry
Installation of carbon-sulfur bonds
Alkylation of sodium sulfide give thioethers:
Na2S + 2 RX → R2S + 2 NaX
Even aryl halides participate in this reaction. By a broadly similar process sodium sulfide can react with alkenes in the thiol-ene reaction to give thioethers.
Sodium sulfide can be used as nucleophile in Sandmeyer type reactions.
Reducing agent
Aqueous solution of sodium sulfide will reduce nitro groups to amine. This conversion is applied to production of some azo dyes since other reducible groups, e.g. azo group, remain intact. The reduction of nitro aromatic compounds to amines using sodium sulfide is known as the Zinin reaction in honor of its discoverer. Hydrated sodium sulfide reduces 1,3-dinitrobenzene derivatives to the 3-nitroanilines.
Other reactions
Sulfide has also been employed in photocatalytic applications.
Safety
Consisting of the equivalent of sodium hydroxide, sodium sulfide is strongly alkaline and can cause chemical burns. It reacts rapidly with acids to produce hydrogen sulfide, which is highly toxic.
| Physical sciences | Sulfide salts | Chemistry |
2918131 | https://en.wikipedia.org/wiki/Rabbitfish | Rabbitfish | Rabbitfishes or spinefoots, genus Siganus, are perciform fishes in the family Siganidae. It is the only extant genus in its family and has 29 species. In some now obsolete classifications, the species having prominent face stripes—colloquially called foxfaces–are in the genus Lo. Other species, such as the masked spinefoot (S. puellus), show a reduced form of the stripe pattern. Rabbitfishes are native to shallow waters in the Indo-Pacific, but S. luridus and S. rivulatus have become established in the eastern Mediterranean via Lessepsian migration. They are commercially important food fish, and can be used in the preparation of dishes such as bagoong.
Taxonomy
The genus Siganus was described in 1775 by the Danish zoologist Johan Christian Fabricius with Siganus rivulatus, a species also described by Fabricius in 1775, designated as the type species. The description was based on notes taken by the naturalist Peter Forsskål when he was on the Danish Arabia expedition (1761–67) and was published in Carsten Niebuhr's Descriptiones animalium avium, amphibiorum, piscium, insectorum, vermium; quae in itinere orientali observavit Petrus Forskål. Post mortem auctoris edidit Carsten Niebuhr. Catalog of Fishes lists the authority as " Fabricius [J. C.] (ex Forsskål) in Niebuhr 1775" and states that the genus is valid as "Siganus Fabricius 1775".
Carl Linnaeus originally described the genus Teuthis, with the type species being Teuthis hepatus. One of the type specimens he used looks like Siganus javus, although the other is definitely not a rabbitfish, and the International Commission on Zoological Nomenclature has been asked to suppress the name Teuthis in favour of Siganus to reflect the prevailing usage.
The name Siganus is a latinisation of the local Arabic name for the marbled rabbitfish (S. rivulatus) in Yemen, Sidjan which can also be written as Sigian, and means "rabbitfish".
In 2007 Kurriwa et al., outlined a way to split the genus—if the scientific community so desires:
An ancient group containing e.g. S. woodlandi
Another fairly small group containing, e.g., the S. canaliculatus/S. fuscescens) complex
The remainder of Siganus, including the foxfaces
Other lineages might exist and make obsolete the somewhat weak distinction between the second and third groups. Also, it is not known where the type species S. rivulatus would fall, hence names for these three subgenera or genera are not established at present.
Hybridizaton has played a role in the evolution of the Siganidae, as evidenced by comparison of mtDNA cytochrome b and nDNA internal transcribed spacer 1 sequence data. Evidence exists of interbreeding between S. guttatus and S. lineatus, as well as between S. doliatus and S. virgatus.
Also, either females of the last common ancestor of S. puellus and the S. punctatus interbred with females ancestral to the main non-foxface lineage, or males of the former hybridized with females of the last common ancestor of S. punctatissimus and the foxfaces, while males of the latter mated with females of the original foxface species.
An individual was found that looked like a slightly aberrant blue-spotted spinefoot (S. corallinus). On investigation, it turned out to be an offspring of a hybrid between a female of that species and a male masked spinefoot, which had successfully backcrossed with the blue-spotted spinefoot.
Species
As noted above, several presumed species are suspected to actively interbreed even today; these might warrant merging as a single species. This applies to the white-spotted spinefoot (S. canaliculatus) and the mottled spinefoot (S. fuscescens), and to the blotched foxface (S. unimaculatus) and the foxface rabbitfish (S. vulpinus). Alternatively they might be very recently evolved species that have not yet undergone complete lineage sorting, but their biogeography suggests that each group is just color morphs of a single species. On the other hand, the morphologically diverse blue-spotted spinefoot (S. corallinus) might represent more than one species; orange individuals are found at the north of its range, while yellow ones occur to the south, and these two may be completely parapatric.
There are currently 29 recognized species in this genus:
Siganus argenteus (Quoy & Gaimard, 1825) (Streamlined spinefoot)
Siganus canaliculatus (M. Park, 1797) (White-spotted spinefoot)
Siganus corallinus (Valenciennes, 1835) (Blue-spotted spinefoot)
Siganus doliatus Guérin-Méneville, 1829 (Barred spinefoot)
Siganus fuscescens (Houttuyn, 1782) (Mottled spinefoot)
Siganus guttatus (Bloch, 1787) (Goldlined spinefoot)
Siganus insomnis Woodland & R. C. Anderson, 2014 (Bronze-lined rabbitfish)
Siganus javus (Linnaeus, 1766) (Streaked spinefoot)
Siganus labyrinthodes (Bleeker, 1853) (Labyrinth spinefoot)
Siganus lineatus (Valenciennes, 1835) (Golden-lined spinefoot)
Siganus luridus (Rüppell, 1829) (Dusky spinefoot)
Siganus magnificus (G. H. Burgess, 1977) (Magnificent rabbitfish)
Siganus niger Woodland, 1990 (Black foxface)
Siganus puelloides Woodland & Randall, 1979 (Blackeye rabbitfish)
Siganus puellus (Schlegel, 1852) (Masked spinefoot)
Siganus punctatissimus Fowler & B. A. Bean, 1929 (Peppered spinefoot)
Siganus punctatus (Schneider & Forster, 1801) (Goldspotted spinefoot)
Siganus randalli Woodland, 1990 (Variegated spinefoot)
Siganus rivulatus Forsskål & Niebuhr, 1775 (Marbled spinefoot)
Siganus spinus (Linnaeus, 1758) (Little spinefoot)
Siganus stellatus (Forsskål, 1775) (Brown-spotted spinefoot)
Siganus sutor (Valenciennes, 1835) (Shoemaker spinefoot)
Siganus trispilos Woodland & G. R. Allen, 1977 (Threeblotched rabbitfish)
Siganus unimaculatus (Evermann & Seale, 1907) (Blotched foxface)
Siganus uspi Gawel & Woodland, 1974 (Bicolored foxface)
Siganus vermiculatus (Valenciennes, 1835) (Vermiculated spinefoot)
Siganus virgatus (Valenciennes, 1835) (Barhead spinefoot)
Siganus vulpinus (Schlegel & J. P. Müller, 1845) (Foxface)
Siganus woodlandi Randall & Kulbicki, 2005
Characteristics
Rabbitfishes have laterally compressed, oval bodies which may be deep, or slender. A few species have a tubular snout. The mouth is very small and is with non protractile jaws which have one row of compressed, closely set, incisor-like teeth in each jaw. The teeth overlap slightly and create a beak like structure. The dorsal fin has 13 robust spines and 10 soft rays and the front spine is short, sharp and points forward, sometimes projecting from its "pocket" but it may be enfolded. The anal fin has 7 robust spines and 9 soft rays. The pelvic fins have 2 spines with 3 soft rays between them; this characteristic is unique to the Siganidae. There is a membrane which extends from the inner pelvic fin spine to the belly with the anus sitting between these membranes. The tiny scales are cycloid and may be absent from the head region. If present on the head they are restricted to a small area of the cheek under the eye. The fin spines are equipped with well-developed venom glands. The sting is very painful, but it is generally not considered medically significant in healthy adults. They range in maximum total lengths of in the case of the blotched foxface (S. unimaculatus) to in the streaked spinefoot (S. javus).
Distribution and habitat
Rabbitfishes are found in the Indo-Pacific from the Red Sea and the coast of eastern Africa through the Pacific Ocean as far as Pitcairn Island. Two Red Sea species S. rivulatus and S. luridus have invaded the Mediterranean Sea through the Suez Canal, a process known as Lessepsian migration. These fishes are found in inshore tropical and subtropical waters where they occur in reefs, lagoons, mangroves and seagrass beds.
Biology
All rabbitfish are diurnal; some live in schools, while others live more solitary lives among the corals. Rabbitfish sleep in crevices in the reef matrix at night. While sleeping, the rabbitfish Siganus canaliculatus was observed being cleaned by the cleaner shrimp Urocaridella antonbruunii. They are herbivorous, feeding on benthic algae in the wild. However, Siganus rivulatus was recently observed feeding on jellyfish (Scyphozoa) and comb jellies (Ctenophora) in the Red Sea. Also Siganus fuscescens have been observed eating prawns and other baits, suggesting that some species are opportunistic omnivorous feeders. The live passage of benthic organisms in the guts of invasive rabbitfish (ichthyochory) was shown to play a major role in the long distance dispersal and bioinvasion of foraminifera. Rabbitfish lay adhesive eggs and some species live as monogamous pairs.
Venom
Rabbitfish have venomous spines in the dorsal and pelvic fins. In at least one species the venom has been found to be similar to that found in stonefish.
Utilization
Rabbitfish can be important species for commercial fisheries, particularly the schooling species. The catch is largely sold fresh but juveniles may be dried or processed to make fish paste. Some species are used in aquaculture and some of the more colorful species are found in the aquarium trade.
Some species have been reported to be hallucinogenic.
| Biology and health sciences | Acanthomorpha | Animals |
949263 | https://en.wikipedia.org/wiki/Stele%20%28biology%29 | Stele (biology) | In a vascular plant, the stele is the central part of the root or stem containing the tissues derived from the procambium. These include vascular tissue, in some cases ground tissue (pith) and a pericycle, which, if present, defines the outermost boundary of the stele. Outside the stele lies the endodermis, which is the innermost cell layer of the cortex.
The concept of the stele was developed in the late 19th century by French botanists P. E. L. van Tieghem and H. Doultion as a model for understanding the relationship between the shoot and root, and for discussing the evolution of vascular plant morphology. Now, at the beginning of the 21st century, plant molecular biologists are coming to understand the genetics and developmental pathways that govern tissue patterns in the stele. Moreover, physiologists are examining how the anatomy (sizes and shapes) of different steles affect the function of organs.
Protostele
The earliest vascular plants had stems with a central core of vascular tissue. This consisted of a cylindrical strand of xylem, surrounded by a region of phloem. Around the vascular tissue there might have been an endodermis that regulated the flow of water into and out of the vascular system. Such an arrangement is termed a protostele.
There are usually three basic types of protostele:
haplostele – consisting of a cylindrical core of xylem surrounded by a ring of phloem. An endodermis generally surrounds the stele. A centrarch (protoxylem in the center of a metaxylem cylinder) haplostele is prevalent in members of the rhyniophyte grade, such as Rhynia.
actinostele – a variation of the protostele in which the core is lobed or fluted. This stele is found in many species of club moss (Lycopodium and related genera). Actinosteles are typically exarch (protoxylem external to the metaxylem) and consist of several to many patches of protoxylem at the tips of the lobes of the metaxylem. Exarch protosteles are a defining characteristic of the lycophyte lineage.
plectostele – a protostele in which plate-like regions of xylem appear in transverse section surrounded by phloem tissue, thus appearing to form alternating bands. These discrete plates are interconnected in longitudinal section. Some modern club mosses have plectosteles in their stems. The plectostele may be derived from the actinostele.
Siphonostele
Siphonosteles have a central region of ground tissue called the pith, with the vascular strand comprising a hollow cylinder surrounding the pith. Siphonosteles often have interruptions in the vascular strand where leaves (typically megaphylls) originate (called leaf gaps).
Siphonosteles can be called ectophloic (phloem present only external to the xylem) or they can be amphiphloic (with phloem both external and internal to the xylem). Among living plants, many ferns and some Asterid flowering plants have an amphiphloic stele.
An amphiphloic siphonostele can be called a solenostele, or this term may be used to refer to cases where the cylinder of vascular tissue contains no more than one leaf gap in any transverse section (i.e. has non-overlapping leaf gaps). This type of stele is primarily found in fern stems today. Where there are large overlapping leaf gaps (so that multiple gaps in the vascular cylinder exist in any one transverse section), the term dictyostele may be used. The numerous leaf gaps and leaf traces give a dictyostele the appearance of many isolated islands of xylem surrounded by phloem. Each of the apparently isolated units of a dictyostele that serve a single leaf can be called a meristele. Among living plants, this type of stele is found only in the stems of ferns.
Most seed plant stems possess a vascular arrangement which has been interpreted as a derived siphonostele, and is called a eustele – in this arrangement, the primary vascular tissue consists of vascular bundles, usually in one or two rings around the pith. In addition to being found in stems, the eustele appears in the roots of monocot flowering plants. The vascular bundles in a eustele can be collateral (with the phloem on only one side of the xylem) or bicollateral (with phloem on both sides of the xylem, as in some Solanaceae).
There is also a variant of the eustele found in monocots like maize and rye. The variation has numerous scattered bundles in the stem and is called an atactostele (characteristic of monocot stems). However, it is really just a variant of the eustele.
| Biology and health sciences | Plant stem | Biology |
949620 | https://en.wikipedia.org/wiki/Placoderm | Placoderm | Placoderms (from Greek πλάξ (plax, plakos) 'plate' and δέρμα (derma) 'skin') are vertebrate animals of the class Placodermi, an extinct group of prehistoric fish known from Paleozoic fossils during the Silurian and the Devonian periods. While their endoskeletons are mainly cartilaginous, their head and thorax were covered by articulated armoured plates (hence the name), and the rest of the body was scaled or naked depending on the species.
Placoderms were among the first jawed fish (their jaws likely evolved from the first pair of gill arches), as well as the first vertebrates to have true teeth. They were also the first fish clade to develop pelvic fins, the second set of paired fins and the homologous precursor to hindlimbs in tetrapods. 380-million-year-old fossils of three other genera, Incisoscutum, Materpiscis and Austroptyctodus, represent the oldest known examples of live birth.
Placoderms are thought to be paraphyletic, consisting of several distinct outgroups or sister taxa to all living jawed vertebrates, which originated among their ranks. In contrast, one 2016 analysis concluded that Placodermi is likely monophyletic.
The first identifiable placoderms appear in the fossil record during the late Llandovery epoch of the early Silurian. They eventually outcompeted the previously dominant marine arthropods (e.g. eurypterids) and cephalopod molluscs (e.g. orthocones), producing some of the first and most infamous vertebrate apex predators such as Eastmanosteus, Dinichthys and the massive Dunkleosteus. Various groups of placoderms were diverse and abundant during the Devonian, but all placoderms became extinct at the end-Devonian Hangenberg event 358.9 million years ago, leaving the niches open for the osteichthyan and chondrichthyan survivors who subsequently radiated during the Carboniferous.
Characteristics
Many placoderms, particularly the Rhenanida, Petalichthyida, Phyllolepida, and Antiarchi, were bottom-dwellers. In particular, the antiarchs, with their highly modified, jointed bony pectoral fins, were highly successful inhabitants of Middle-Late Devonian freshwater and shallow marine habitats, with the Middle to Late Devonian genus, Bothriolepis, known from over 100 valid species. The vast majority of placoderms were predators, many of which lived at or near the substrate. Many, primarily the arthrodires, were active, nektonic predators that dwelled in the middle to upper portions of the water column. A study of the arthrodire Compagopiscis published in 2012 concluded that placoderms (at least this particular genus) likely possessed true teeth contrary to some early studies. The teeth had well defined pulp cavities and were made of both bone and dentine. However, the tooth and jaw development were not as closely integrated as in modern gnathostomes. These teeth were likely homologous to the teeth of other gnathostomes.
One of the largest known arthrodires, Dunkleosteus terrelli, was long, and is presumed to have had a large distribution, as its remains have been found in Europe, North America and possibly Morocco. Some paleontologists regard it as the world's first vertebrate "superpredator", preying upon other predators. Other, smaller arthrodires, such as Fallacosteus and Rolfosteus, both of the Gogo Formation of Western Australia, had streamlined, bullet-shaped head armor, and Amazichthys, with morphology like that of other fast-swimming pelagic organisms, strongly supporting the idea that many, if not most, arthrodires were active swimmers, rather than passive ambush-hunters whose armor practically anchored them to the sea floor. Some placoderms were herbivorous, such as the Middle to Late Devonian arthrodire Holonema, and some were planktivores, such as the gigantic arthrodire Titanichthys, various members of Homostiidae, and Heterosteus.
Extraordinary evidence of internal fertilization in a placoderm was afforded by the discovery in the Gogo Formation, near Fitzroy Crossing, Kimberley, Western Australia, of a small female placoderm, about in length, which died in the process of giving birth to a 6 cm ( in) offspring and was fossilized with the umbilical cord intact. The fossil, named Materpiscis attenboroughi (after scientist David Attenborough), had eggs which were fertilized internally, the mother providing nourishment to the embryo and giving birth to live young. With this discovery, the placoderm became the oldest vertebrate known to have given birth to live young ("viviparous"), pushing the date of first viviparity back some 200 million years earlier than had been previously known. Specimens of the arthrodire Incisoscutum ritchei, also from the Gogo Formation, have been found with embryos inside them indicating this group also had live bearing ability. The males reproduced by inserting a long clasper into the female. Elongated basipterygia are also found on the phyllolepid placoderms, such as Austrophyllolepis and Cowralepis, both from the Middle Devonian of Australia, suggesting that the basipterygia were used in copulation.
The placoderm claspers are not homologous with the claspers in cartilaginous fishes. The similarities between the structures has been revealed to be an example of convergent evolution. While the claspers in cartilaginous fishes are specialized parts of their paired pelvic fins that have been modified for copulation due to changes in the hox genes hoxd13, the origin of the mating organs in placoderms most likely relied on different sets of hox genes and were structures that developed further down the body as an extra and independent pair of appendages, but which during development turned into body parts used for reproduction only. Because they were not attached to the pelvic fins, as are the claspers in fish like sharks, they were much more flexible and could probably be rotated forward.
A study on Kolymaspis showcases that the vertebrate shoulder girdle evolved from gill arches.
Evolution and extinction
It was thought for a time that placoderms became extinct due to competition from the first bony fish and early sharks, given a combination of the supposed inherent superiority of bony fish and the presumed sluggishness of placoderms. With more accurate summaries of prehistoric organisms, it is now thought that they systematically died out as marine and freshwater ecologies suffered from the environmental catastrophes of the Late Devonian and end-Devonian extinctions.
Fossil record
The earliest identifiable placoderm fossils are of Chinese origin and date to the early Silurian. At that time, they were already differentiated into antiarchs and arthrodires, as well as other, more primitive, groups. Earlier fossils of basal placoderms have not yet been discovered.
The Silurian fossil record of the placoderms is both literally and figuratively fragmented. Until the discovery of Silurolepis (and then, the discoveries of Entelognathus and Qilinyu), Silurian-aged placoderm specimens consisted of fragments. Some of them have been tentatively identified as antiarch or arthrodire due to histological similarities; and many of them have not yet been formally described or even named. The most commonly cited example of a Silurian placoderm, Wangolepis of Silurian China and possibly Vietnam, is known only from a few fragments that currently defy attempts to place them in any of the recognized placoderm orders. So far, only three officially described Silurian placoderms are known from more than scraps:
the basal antiarch Silurolepis, from the Ludlow epoch of Yunnan, China, known from an almost complete thoracic armor
Entelognathus, a placoderm incertae sedis that combines features of primitive arthrodires with jaw anatomy otherwise only seen in bony fish and tetrapods.
Qilinyu, a close relative of Entelognathus that further links Entelognathus as a transitional form between placoderms and other stem-gnathostomes and crown-group gnathostomes.
The first officially described Silurian placoderm is an antiarch, Shimenolepis, which is known from distinctively ornamented plates from Hunan, China. It was originally considered to be from the late Llandovery, although later study reconsidered its age at Ludfordian. Shimenolepis plates are very similar to the early Devonian yunnanolepid Zhanjilepis, also known from distinctively ornamented plates. In 2022, Xiushanosteus is described from complete fossils from Telychian, late Llandovery of Chongqing, China.
Paleontologists and placoderm specialists suspect that the scarcity of placoderms in the Silurian fossil record is due to placoderms' living in environments unconducive to fossil preservation, rather than a genuine scarcity. This hypothesis helps to explain the placoderms' seemingly instantaneous appearance and diversity at the very beginning of the Devonian.
During the Devonian, placoderms went on to inhabit and dominate almost all known aquatic ecosystems, both freshwater and saltwater. But this diversity ultimately suffered many casualties during the extinction event at the Frasnian–Famennian boundary, the Late Devonian extinctions. The remaining species then died out during the end-Devonian extinction; not a single placoderm species has been confirmed to have survived into the Carboniferous.
History of study
The earliest studies of placoderms were published by Louis Agassiz, in his five volumes on fossil fishes, 1833–1843. In those days, placoderms were thought to be shelled jawless fish akin to ostracoderms. Some naturalists even suggested that they were shelled invertebrates or even turtle-like vertebrates.
In the late 1920s, Dr. Erik Stensiö, at the Swedish Museum of Natural History in Stockholm, established the details of placoderm anatomy and identified them as true jawed fishes related to sharks. He took fossil specimens with well-preserved skulls and ground them away, one tenth of a millimeter at a time. After each layer had been removed, he made an imprint of the next surface in wax. Once the specimens had been completely ground away (and so destroyed), he made enlarged, three-dimensional models of the skulls to examine the anatomical details more thoroughly. Many other placoderm specialists thought that Stensiö was trying to shoehorn placoderms into a relationship with sharks; however, as more fossils were found, placoderms were accepted as a sister group of chondrichthyans.
Much later, the exquisitely preserved placoderm fossils from Gogo reef changed the picture again. They showed that placoderms shared anatomical features not only with chondrichthyans but with other gnathostome groups as well. For example, Gogo placoderms show separate bones for the nasal capsules as in gnathostomes; in both sharks and bony fish those bones are incorporated into the braincase.
Placoderms also share certain anatomical features only with the jawless osteostracans; because of this, the theory that placoderms are the sister group of chondrichthyans has been replaced by the theory that placoderms are a group of basal gnathostomes.
Taxonomy and phylogeny
Currently, Placodermi are divided into eight recognized orders. There are two further controversial orders: One is the monotypic Stensioellida, containing the enigmatic Stensioella; the other is the equally enigmatic Pseudopetalichthyida. These orders are considered to be basal or primitive groups within Placodermi, though their precise placement within the class remains unsure. Fossils of both are currently known only from the Hunsruck lagerstatten.
Placoderm orders
Arthrodira
Arthrodira ("jointed neck") were the most diverse and numerically successful of the placoderm orders, occupying roles from giant apex predators to detritus-nibbling bottom dwellers. They had a movable joint between armour surrounding the head and body. As the lower jaw moved down, the head shield moved, allowing for a larger opening. All arthrodires, save for Compagopiscis, lacked teeth, and used instead the sharpened edges of a bony plate, termed a "tooth plate," as a biting surface (Compagopiscis had true teeth in addition to tooth plates). The eye sockets are protected by a bony ring, a feature shared by birds and some ichthyosaurs. Early arthrodires, such as the genus Arctolepis, were well-armoured fishes with flattened bodies. The largest member of this group, Dunkleosteus, was a true "superpredator" of the latest Devonian period, reaching 3 to as much as 8 metres in length. In contrast, the long-nosed Rolfosteus measured just 15 cm. Fossils of Incisoscutum have been found containing unborn fetuses, indicating that arthrodires gave birth to live young.
Antiarchi
Antiarchi ("opposite anus") were the second most successful order of placoderms known, after the Arthrodira. The order's name was coined by Edward Drinker Cope, who, after incorrectly identifying the first fossils as being those of an armored tunicate, mistakenly thought the eye-hole was the mouth, and the opening for the anal siphon was on the other side of the body, as opposed to having both oral and anal siphons together at one end. The front portions of their bodies were heavily armoured, to the point of literally resembling a box with eyes, with the sometimes scaled, sometimes naked rear portions often becoming sinuous, particularly with later forms. The pair of pectoral fins were modified into a pair of caliper-like, or arthropod-like limbs. In primitive forms, such as Yunnanolepis, the limbs were thick and short, while in advanced forms, such as Bothriolepis, the limbs were long and had elbow-like joints. The function of the limbs is still not perfectly understood, but most hypothesize that they helped their owners pull themselves across the substrate, as well as allowing their owners to bury themselves into the substrate.
Brindabellaspida
Brindabellaspis ("Brindabella's shield") was a long-snouted placoderm from the Early Devonian. When it was first discovered in 1980, it was originally regarded as a weejasperaspid acanthothoracid due to anatomical similarities with the other species found at the same locality. According to Philippe Janvier, anatomical similarities in the brain of Brindabellaspis stensioi and the brain of a jawless fish suggest it is a basal placoderm closest to the ancestral placoderm. Various Early to Middle Devonian placoderm incertae sedis have also been inserted in the order.
Phyllolepida
Phyllolepida ("leaf scales") were flattened placoderms found throughout the world. Like other flattened placoderms they were bottom-dwelling predators that ambushed prey. Unlike other flattened placoderms, they were freshwater fish. Their armour was made of whole plates, rather than the numerous tubercles and scales of Petalichthyida. The eyes were on the sides of the head, unlike visual bottom-dwelling predators, such as stargazers or flatfish, which have eyes on the top of their head. The orbits for the eyes were extremely small, suggesting the eyes were vestigial and that the phyllolepids may have been blind.
Ptyctodontida
Ptyctodontida ("folded teeth") were lightly armoured placoderms with big heads, big eyes and long bodies. They have a strong but superficial resemblance to modern day chimaeras. Their armour was reduced to a pattern of small plates around the head and neck. Like the extinct and related acanthothoracids, and the living and unrelated holocephalians, most of the ptyctodontids are thought to have lived near the sea bottom and preyed on shellfish. On account of their lack of armour, some paleontologists have suggested that the Ptyctodontida were not placoderms, but holocephalians or the ancestors of holocephalians. Anatomical examinations of whole fossil specimens have shown that the similarities between these two groups are superficial. The major differences were that holocephalians have shagreen on their skin, while ptyctodontids do not; the armoured plates and scales of holocephalians are made of dentine, while those of ptyctodontids are made of bone; the craniums of holocephalians are similar to sharks, while those of ptyctodontids are similar to those of other placoderms; and, most importantly, that holocephalians have true teeth, while ptyctodonts have beak-like tooth plates. Ptyctodontids were sexually dimorphic, with the males having pelvic claspers and possibly claspers on the head as well.
Rhenanida
Rhenanida ("Rhine fish") were flattened, ray-like, bottom-dwelling predators with large, upturned mouths that lived in marine environments. The rhenanids were once presumed to be the most primitive, or at least the closest to the ancestral placoderm, as their armour was made of unfused components—a mosaic of tubercles—as opposed to the solidified plates of "advanced" placoderms, such as antiarchs and arthrodires. However, through comparisons of skull anatomies, rhenanids are now considered to be the sister group of the antiarchs. When rhenanids die, their "mosaics" come apart, and it has been suggested that the rarity of rhenanids in the fossil record reflects postmortem disassociation, and is not an actual rarity of the species.
Acanthothoraci
Acanthothoraci ("spine chests") were a group of chimaera-like placoderms closely related to the rhenanid placoderms. Superficially, acanthoracids resembled scaly chimaeras or small, scaly arthrodires with blunt rostrums. They were distinguished from chimaeras by a pair of large spines that emanate from their chests, the presence of large scales and plates, tooth-like beak plates, and the typical bone-enhanced placoderm eyeball. They were distinguished from other placoderms due to differences in the anatomy of their skulls, and due to patterns on the skull plates and thoracic plates that are unique to this order. From what can be inferred from the mouthplates of fossil specimens, acanthothoracids were shellfish hunters ecologically similar to modern-day chimaeras. Competition with their relatives, the ptyctodont placoderms, may have been one of the main reasons for the acanthothoracids' extinction prior to the mid-Devonian extinction event.
Petalichthyida
Petalichthyida ("thin-plated fish") were small, flattened placoderms, typified by their splayed fins and numerous tubercles that decorated all of the plates and scales of their armour. They reached a peak in diversity during the Early Devonian and were found throughout the world. The petalichthids Lunaspis and Wijdeaspis are among the best known. There was an independent diversification event that occurred in what is now Southern China, producing a handful of unique genera that were once placed in their own order, "Quasipetalichthyida", named after the first discovered species there, Quasipetalichthys haikouensis. Soon after the petalichthids' diversification, they went into decline. Because they had compressed body forms, it is supposed they were bottom-dwellers that pursued or ambushed smaller fish. Their diet is not clear, as none of the fossil specimens found have preserved mouth parts.
Pseudopetalichthyida
Pseudopetalichthyida ("false petalichthyids") is a group of elongated, possibly flattened fishes comprising three, poorly preserved and poorly studied genera. It is known only from rare fossils in Lower Devonian strata in Hunsrück, Germany. Like Stensioella heintzi, and the Rhenanida, the pseudopetalichthids had armour made up of a mosaic of tubercles. Like Stensioella heintzi, the pseudopetalichthids' placement within Placodermi is suspect. The matter is not easy to resolve because there are no complete, undamaged and articulated specimens. The anatomical studies done on the crushed specimens that have been found indicate that if they are placoderms, they may be a group more advanced than the ptyctodonts. As such, placoderm experts consider Pseudopetalichthyida to be the sister group of the Arthrodires + Phyllolepida + Antiarchi trichotomy and the Acanthothoraci + Rhenanida dichotomy.
Stensioellida
Stensioellida ("[Heintz's] little Stensio") contains another problematic placoderm of uncertain affinity, known only from the Lower Devonian Hunsrück slates of Germany. Stensioella was a thin fish that, when alive, looked vaguely like an elongated ratfish, or a skinny Gemuendina with thin, strap-like pectoral fins. Similar to those of the Rhenanida, its armour was a complex mosaic of small, scale-like tubercles. The shoulder joints of its armour are similar to other placoderms, and there are superficial similarities in skull plates, and even more superficial similarities between its tubercles and the tubercles of the rhenanids. It is tentatively placed within Placodermi as a primitive placoderm, though some paleontologists believe the rationale for the placement is inadequate. The paleontologist Philippe Janvier, as well as other paleontologists, has suggested that Stensioella is not a placoderm, but instead is a holocephalian. If this is true, then the holocephalians diverged from sharks before the Chondrichthyan Devonian radiation. Critics of Janvier's position say that aside from a bodyplan superficially similar to primitive holocephalians, the two groups have little else in common anatomically.
Cladogram
The following cladogram shows the interrelationships of placoderms according to Carr et al. (2009):
However, the cladogram had changed significantly over the years, and the placoderms are now thought to be paraphyletic, with some being more closer to the Eugnathostomata than others. The updated cladogram (Zhu et al., 2016):
| Biology and health sciences | Fishes | null |
950991 | https://en.wikipedia.org/wiki/Little%20finger | Little finger | The little finger or pinkie, also known as the baby finger, fifth digit, or pinky finger, is the most ulnar and smallest digit of the human hand, and next to the ring finger.
Etymology
The word "pinkie" is derived from the Dutch word pink, meaning "little finger".
The earliest recorded use of the term "pinkie" is from Scotland in 1808. The term (sometimes spelled "pinky") is common in Scottish English and American English, and is also used extensively in other Commonwealth countries such as New Zealand, Canada, and Australia.
Nerves and muscles
There are nine muscles that control the fifth digit:
Three in the hypothenar eminence, two extrinsic flexors, two extrinsic extensors, and two more intrinsic muscles:
Hypothenar eminence:
Opponens digiti minimi muscle
Abductor minimi digiti muscle (adduction from third palmar interossei)
Flexor digiti minimi brevis (the "longus" is absent in most humans)
Two extrinsic flexors:
Flexor digitorum superficialis
Flexor digitorum profundus
Two extrinsic extensors:
Extensor digiti minimi muscle
Extensor digitorum
Two intrinsic hand muscles:
Fourth lumbrical muscle
Third palmar interosseous muscle
Note: the dorsal interossei of the hand muscles do not have an attachment to the fifth digit
Cultural significance
Gestures
Among American children, a "pinky swear" or "pinky promise" is made when a person wraps one of their pinky fingers around another person's pinky and makes a promise. Something similar is also seen in China and Korea, where people link their pinky fingers and then stamp their thumbs together to make a yaksok (promise).
Other gestures where the little finger is extended include the shaka sign and sign of the horns.
Among members of the Japanese yakuza (gangsters), the penalty for various offenses is removal of parts of the little finger (known as yubitsume).
It is a common joke that one should extend their little finger when drinking from a teacup in imitation of a passé upper-class tradition. This practice is generally deprecated by etiquette guides as a sign of snobbery amongst the socially inferior, with various cultural theories as to the origin of the practice including the idea that finger food should be eaten with only the first three digits.
The messaging application Teams from Microsoft has an emoji which is a representation of a closed hand with the little finger raised. The description is "Nature's call" which is a polite euphemism used when someone feels a need to urinate or defecate.
Rings
The signet ring is traditionally worn on the little finger of a gentleman's left hand, a practice still common especially in the United Kingdom, Australia, and European cultures. A signet ring is considered part of the regalia of many European monarchies, and also of the Pope, with the ring always worn on the left little finger. In modern times the location of the signet ring has relaxed, with examples worn on various different digits, although little fingers still tend to be the most usual.
The Iron Ring is a symbolic ring worn by most Canadian engineers. The Ring is a symbol of both pride and humility for the engineering profession, and is always worn on the little finger of the dominant hand. In the United States, the Engineer's Ring is a stainless steel ring worn on the fifth digit of the working hand by engineers who belong to the Order of the Engineer and have accepted the Obligation of an Engineer.
Utility
The little finger is often used as a support when smartphone users type one-handed. The little finger is positioned underneath the phone, allowing it to be propped with the three middle fingers, and the user to type with their thumb.
Some users reported dents on their little finger and pain in the hand after prolonged use in this way, doctors referred to this as "iPhone pinky" or "smartphone pinky". The skin indentations were reported to be nothing of alarm, as they disappeared on their own after a short while without cell phone use.
| Biology and health sciences | Human anatomy | Health |
950992 | https://en.wikipedia.org/wiki/Elevator%20%28aeronautics%29 | Elevator (aeronautics) | Elevators are flight control surfaces, usually at the rear of an aircraft, which control the aircraft's pitch, and therefore the angle of attack and the lift of the wing. The elevators are usually hinged to the tailplane or horizontal stabilizer. They may be the only pitch control surface present, and are sometimes located at the front of the aircraft (early airplanes and canards) or integrated into a rear "all-moving tailplane", also called a slab elevator or stabilator.
Elevator control effectiveness
The elevator is a usable up and down system that controls the plane, horizontal stabilizer usually creates a downward force which balances the nose down moment created by the wing lift force, which typically applies at a point (the wing center of lift) situated aft of the airplane's center of gravity. The effects of drag and changing the engine thrust may also result in pitch moments that need to be compensated with the horizontal stabilizer.
Both the horizontal stabilizer and the elevator contribute to pitch stability, but only the elevators provide pitch control. They do so by decreasing or increasing the downward force created by the stabilizer:
an increased downward force, produced by up elevator, forces the tail down and the nose up. At constant speed, the wing's increased angle of attack causes a greater lift to be produced by the wing, accelerating the aircraft upwards. The drag and power demand also increase;
a decreased downward force at the tail, produced by down elevator, causes the tail to rise and the nose to lower. At constant speed, the decrease in angle of attack reduces the lift, accelerating the aircraft downwards.
On many low-speed aircraft, a trim tab is present at the rear of the elevator, which the pilot can adjust to eliminate forces on the control column at the desired attitude and airspeed. Supersonic aircraft usually have all-moving tailplanes (stabilators), because shock waves generated on the horizontal stabilizer greatly reduce the effectiveness of hinged elevators during supersonic flight. Delta winged aircraft combine ailerons and elevators –and their respective control inputs– into one control surface called an elevon.
Elevators' location
Elevators are usually part of the tail, at the rear of an aircraft. In some aircraft, pitch-control surfaces are in the front, ahead of the wing. In a two-surface aircraft this type of configuration is called a canard (the French word for duck) or a tandem wing. The Wright Brothers' early aircraft were of the canard type; Mignet Pou-du-Ciel and Rutan Quickie are of tandem type. Some early three surface aircraft had front elevators (Curtiss/AEA June Bug); modern three surface aircraft may have both front (canard) and rear elevators (Grumman X-29).
Research
Several technology research and development efforts exist to integrate the functions of aircraft flight control systems such as ailerons, elevators, elevons, flaps and flaperons into wings to perform the aerodynamic purpose with the advantages of less: mass, cost, drag, inertia (for faster, stronger control response), complexity (mechanically simpler, fewer moving parts or surfaces, less maintenance), and radar cross section for stealth. These may be used in many unmanned aerial vehicles (UAVs) and 6th generation fighter aircraft. Two promising approaches are flexible wings, and fluidics.
In flexible wings, much or all of a wing surface can change shape in flight to deflect air flow. The X-53 Active Aeroelastic Wing is a NASA effort. The Adaptive Compliant Wing is a military and commercial effort.
In fluidics, forces in vehicles occur via circulation control, in which larger more complex mechanical parts are replaced by smaller simpler fluidic systems (slots which emit air flows) where larger forces in fluids are diverted by smaller jets or flows of fluid intermittently, to change the direction of vehicles. In this use, fluidics promises lower mass, costs (up to 50% less), and very low inertia and response times, and simplicity.
Gallery
| Technology | Aircraft components | null |
952459 | https://en.wikipedia.org/wiki/Bed%20bug | Bed bug | Bed bugs are parasitic insects from the genus Cimex, who are micropredators that feed on blood, usually at night. Their bites can result in a number of health impacts, including skin rashes, psychological effects, and allergic symptoms. Bed bug bites may lead to skin changes ranging from small areas of redness to prominent blisters. Symptoms may take between minutes to days to appear and itchiness is generally present. Some individuals may feel tired or have a fever. Typically, uncovered areas of the body are affected. Their bites are not known to transmit any infectious disease. Complications may rarely include areas of dead skin or vasculitis.
Bed bug bites are caused primarily by two species of insects: Cimex lectularius (the common bed bug) and Cimex hemipterus, found primarily in the tropics. Their size ranges between 1 and 7 mm. They spread by crawling between nearby locations or by being carried within personal items. Infestation is rarely due to a lack of hygiene but is more common in high-density areas. Diagnosis involves both finding the bugs and the occurrence of compatible symptoms. Bed bugs spend much of their time in dark, hidden locations like mattress seams, or cracks in a wall.
Treatment is directed towards the symptoms. Eliminating bed bugs from the home is often difficult, partly because bed bugs can survive up to approximately 300 days without feeding. Repeated treatments of a home may be required. These treatments may include heating the room to for more than 90 minutes, frequent vacuuming, washing clothing at high temperatures, and the use of various pesticides.
Fossils found in Egypt show bed bugs have been known as human parasites for at least 3,500 years. Despite being nearly eradicated in developed countries after World War II, infestations have increased since the 1990s and bed bugs are now relatively common in all regions of the globe. Experts point to several factors that have contributed to the explosion in infestations over the last three decades: increased immigration and international travel; expanded markets for second-hand goods; a greater focus on control of other pests; the banning of certain pesticides and increased resistance to pesticides still in use.
Effects on humans
Bed bugs infest dwellings and bite people, causing irritation and sometimes other issues. There is no evidence that bed bugs transmit infectious diseases even though they appear physically capable of carrying pathogens and this possibility has been investigated.
Bites
The most common skin findings associated with bed bug bites are itching, flat and bumpy, reddish lesions. Each lesion is about but may be as large as in diameter and there may or may not be a central spot (punctum). Bites are usually present on areas of exposed skin, especially exposed areas not covered by sheets or blankets, such as arms, legs, feet, face or neck. Individual responses to bites vary, ranging from no visible effect (in about 20–70%), to small flat (macular) spots, to the formation of prominent blisters (wheals and bullae) along with intense itching that may last several days. Vesicles and nodules may also form. The lesions due to bites may become secondarily infected due to scratching but systemic effects from bed bug bites are very rare. A central spot of bleeding may also occur due to the release of blood thinning substances in the bug's saliva.
Symptoms may not appear until some days after the bites have occurred. Reactions often become brisker after multiple bites due to possible sensitization to the salivary proteins of the bed bug. Numerous bites may lead to a red rash or hives.
Bedbug bites may cause other symptoms and health issues. Serious allergic reactions including anaphylaxis from the injection of serum and other non-specific proteins have been documented, though rarely. As each bite takes a tiny amount of blood, chronic or severe infestation may lead to anemia. Scratching bites may lead to bacterial skin infection. Systemic poisoning may occur if the bites are numerous. The bite itself may be painful thus resulting in poor sleep and worse work performance.
Bed bugs can feed on warm-blooded animals other than humans, such as pets. The signs left by the bites are the same as in the case of people and cause identical symptoms (skin irritation, scratching etc.). Bed bugs can infest poultry sheds and cause anemia and a decrease in egg production in hens.
Treatment
Treatment of bed bug bites requires keeping the person from being repeatedly bitten, and possible symptomatic use of antihistamines and corticosteroids (either topically or systemically). There however is no evidence that medications improve outcomes, and symptoms usually resolve without treatment in 1–2 weeks.
Other effects of infestation
It is possible that exposure to bed bugs may trigger an asthma attack via the effects of airborne allergens, although evidence of this association is limited.
Serious infestations and chronic attacks can cause anxiety, stress, and sleep difficulties. Development of refractory delusional parasitosis is possible, as a person develops an overwhelming obsession with bed bugs.
Description
Bed bug infestations are primarily the result of two species of insects from genus Cimex: Cimex lectularius (the common bed bug) and Cimex hemipterus (the tropical bed bug). These insects feed exclusively on blood and, at any stage of development, may survive up to 70 days without feeding. Adult Cimex are light brown to reddish-brown, flat, oval, and have no hind wings. The front wings are vestigial and reduced to pad-like structures. Adults grow to long and wide. Female common bed bugs can lay 1–10 eggs per day and 200–500 eggs in their lifetime, whereas female tropical bed bugs can lay about 50 eggs in their lifetime.
Bed bugs have five immature nymph life stages and a final sexually mature adult stage. Bed bugs need at least one blood meal in order to advance to the next stage of development. They shed their skins through ecdysis at each stage, discarding their outer exoskeleton. Newly hatched nymphs are translucent, lighter in color, and become browner as they moult and reach maturity. Bed bugs may be mistaken for other insects, such as booklice, small cockroaches, or carpet beetles; however, when warm and active, their movements are more ant-like, and like most other true bugs, they emit a characteristic disagreeable odor when crushed.
Bed bugs are obligatory bloodsuckers. They have mouth parts that saw through the skin and inject saliva with anticoagulants and painkillers. Sensitivity of humans varies from extreme allergic reaction to no reaction at all (about 20%). The bite usually produces a swelling with no red spot, but when many bugs feed on a small area, reddish spots may appear after the swelling subsides. Bedbugs prefer exposed skin, preferably the face, neck, and arms of a sleeping person.
Bed bugs are attracted to their hosts primarily by carbon dioxide, secondarily by warmth, and also by certain chemicals. There is strong evidence that bed bugs can respond and orient towards human odors, independently of all other host cues. Cimex lectularius feeds only every five to seven days, which suggests that it does not spend the majority of its life searching for a host. When a bed bug is starved, it leaves its shelter and searches for a host. It returns to its shelter after successful feeding or if it encounters exposure to light. Cimex lectularius aggregate under all life stages and mating conditions. Bed bugs may choose to aggregate because of predation, resistance to desiccation, and more opportunities to find a mate. Airborne pheromones are responsible for aggregations.
Infestation
Infestation is rarely caused by a lack of hygiene. Transfer to new places is usually in the personal items of the human they feed upon. Dwellings can become infested with bed bugs in a variety of ways, such as:
Bugs and eggs inadvertently brought in from other infested dwellings on a visiting person's clothing or luggage;
Infested items (such as furniture especially beds or couches, clothing, or backpacks) brought into a home or business;
Proximity of infested dwellings or items, if easy routes are available for travel, e.g. through ducts or false ceilings;
Wild animals (such as bats or birds) that may also harbour bed bugs or related species such as the bat bug;
People visiting an infested area (e.g. dwelling, means of transport, entertainment venue, or lodging) and carrying the bugs to another area on their clothing, luggage, or bodies. Bedbugs are increasingly found in air travel.
Though bed bugs will opportunistically feed on pets, they do not live or travel on the skin of their hosts, and pets are not believed to be a factor in their spread.
Detection
Knowing that symptoms are caused by bedbug bites rather than other causes requires seeking and finding the insect in the sleeping environment, as symptoms are not specific to bedbug bites. Bites by other arthropods cause similar symptoms, even the linear pattern of bites known colloquially as "breakfast, lunch and dinner bites".
Bed bugs can occur singly, but tend to congregate once established. Although strictly parasitic, they spend only a tiny fraction of their lives physically attached to hosts. Once a bed bug finishes feeding, it follows a chemical trail to return to a nearby harborage, commonly in or near beds or couches, where they live in clusters of adults, juveniles, and eggs. These places may include luggage, vehicle interiors, furniture, bedside clutter—even inside electrical sockets or laptop computers. Bed bugs may also lodge near animals that have nested within a dwelling, such as bats, birds, or rodents. They can also survive by feeding on domestic cats and dogs, though humans are the preferred host of C. lectularius.
A severe bedbug infestation can be detected by their characteristic pungent sweet smell, which has been described as like rotting raspberries. Bed bug detection dogs are trained to pinpoint infestations, with a possible accuracy rate between 11% and 83%.
Homemade detectors have been developed. Bedbug detectors, often referred to as "monitors", "traps" or "interceptors", use the lactic acid or carbon dioxide associated with the presence of a human body, or pheromones, to attract and trap bugs in a container. Bedbug detectors can confirm an infestation, but do not trap enough for eradication.
Differential detection
Other conditions which produce symptoms similar to bedbug bites include scabies, gamasoidosis, allergic reactions, mosquito bites, spider bites, flea bites (pulicosis), chicken pox, and bacterial skin infections.
Prevention
To prevent bringing home bed bugs from outside the home, people are advised to take precautions after visiting an infested site or traveling on means of transport that may be infested; precautions include checking shoes on leaving the site, changing clothes outside the house before entering, and putting the used clothes in a clothes dryer outside the house. When visiting a new lodging, it is advised to check the bed before taking suitcases into the sleeping area, and putting the suitcase on a raised stand to make bedbugs less likely to crawl in. Clothes should be hung up or left in the suitcase rather than left on the floor. Additional preventative measures include sealing cracks and crevices (where bed bugs often hide), inspecting furniture, and decontaminating clothes and luggage upon returning home. The founder of a company dedicated to bedbug extermination said that 5% of hotel rooms he booked into were infested. He advised people never to sit down on public transport; check office chairs, plane seats, and hotel mattresses; and monitor and vacuum home beds once a month. Close all wall openings or gaps; bed bugs tend to hide in dark places such as cracks in walls. Second-hand furnishings may harbour bedbugs.
Management
Avoiding repeated bites can be difficult since it usually requires eradicating bed bugs from a home or workplace; eradication is most effective using non-chemical control methods. Non-chemical control methods include vacuuming carpet and furniture (often with scraping) into a disposable bag which is then sealed into a plastic bag to prevent re-infestation. Other methods include removing textile materials from an area and washing them in hot water (at least 60 degrees Celsius) or freezing them at . Most consumer-grade freezers are inadequate to kill bedbugs because they cannot create sufficiently low temperatures. Unremovable textiles such as mattresses can be steamed to at least and this method can penetrate deep into the textile to effectively kill bed bugs in, potentially, under one minute. Heating tents or chambers can be used for infested materials or entire rooms can be heated to at least to effectively eradicate infestation.
There is no evidence to indicate that a combination of non-chemical methods plus insecticides is more effective than non-chemical methods alone with regards to eradication of bed bug infestations.
Insecticides are mostly ineffective for the eradication of bedbug infestations as most bedbugs are resistant to insecticides, including pyrethroids which are found in approximately 90% of commercial grade insecticides. Furthermore, insect foggers (known as "bug bombs") are ineffective in the eradication of bed bug infestation as they are unable to penetrate bed bug harborages. Resistance to pesticides has increased significantly over time, and there are concerns about harm to health from their use.
Once established, bed bugs are extremely difficult to get rid of, particularly in buildings with multiple dwellings, as they may be present in other parts of the building than the dwelling being treated, and can re-establish populations by moving from infested to decontaminated areas.
Mechanical approaches, such as vacuuming up the insects and heat-treating or wrapping mattresses, are effective. An hour at a temperature of or over, or two hours at less than kills them. This may include a domestic clothes drier for fabric or a commercial steamer. Bed bugs and their eggs will die on contact when exposed to surface temperatures above and a steamer can reach well above . A study found 100% mortality rates for bed bugs exposed to temperatures greater than for more than 2 minutes. The study recommended maintaining temperatures of above for more than 20 min to effectively kill all life stages of bed bugs, and because in practice treatment times of 6 to 8 hours are used to account for cracks and indoor clutter. This method is expensive and has caused fires. Starving bedbugs is not effective, as they can survive without eating for 135 to 300 days, depending on temperature.
After the withdrawal of most organochlorine insecticides, that no truly effective insecticides were available. Insecticides that have historically been found effective include pyrethroids, dichlorvos, and malathion. Resistance to pesticides has increased significantly in recent decades. The carbamate insecticide propoxur is highly toxic to bed bugs, but it has potential toxicity to children exposed to it, and the US Environmental Protection Agency has been reluctant to approve it for indoor use. Boric acid, sometimes applied as a safe indoor insecticide against pests such as cockroaches and termites, is not effective against bed bugs because they do not groom.
Distribution
Bed bugs are found everywhere in the world. Before the 1950s about 30% of houses in the United States had bedbugs; this percentage has fallen, which is believed to be partly due to the use of DDT to kill cockroaches. The invention of the vacuum cleaner and simplification of furniture design may have also played a role in the decrease. Others believe it might simply be the cyclical nature of the organism.
However, rates of infestation in developed countries have increased dramatically since the 1980s. This is thought to be due to greater foreign travel; increased immigration from the developing world to the developed world; more frequent exchange of second-hand furnishings among homes; a greater focus on control of other pests, resulting in neglect of bed bug countermeasures; and the banning of effective pesticides coupled with increased resistance to those pesticides still permitted. The decrease in cockroach populations due to insecticide use may have aided bed bugs' resurgence, since cockroaches eat bedbugs. Increasing resistance to DDT and other potent pesticides may have also contributed, along with bans on DDT.
The U.S. National Pest Management Association reported a 71% increase in bed bug calls between 2000 and 2005. The number of reported incidents in New York City alone rose from 500 in 2004 to 10,000 in 2009. In 2013 Chicago was listed as the US city with most bedbug infestation. In response the Chicago City Council passed a bed bug control ordinance to limit spread. Additionally, bed bugs are reaching places in which they never established before, such as southern South America.
The rise in infestations has been hard to track because bed bug infestation is not an easily identifiable problem, and also people do not talk about it. Most reports have been collected from pest-control companies, local authorities, and hotel chains, and the problem may be more severe than is currently believed from reports.
Species
The common bed bug (Cimex lectularius) is the species best adapted to human environments but is also known from birds, Chiroptera, Gallus (chickens and relatives), Myotis myotis, and sheep (Ovis aries). It is found in temperate climates throughout the world. Other species include C. hemipterus, found in tropical regions, which also infests poultry (including Gallus) and bats, and Leptocimex boueti, a relative of C. lectularius adapted for the tropics of West Africa and South America, which infests bats and humans. C. pilosellus and C. pipistrella primarily infest bats, while Haematosiphon inodora, a species of North America, primarily infests poultry.
Evolution
Cimicidae, the ancestor of modern bed bugs, first emerged approximately 115 million years ago, more than 55 million years before bats—their previously presumed initial host—first appeared. From unknown ancestral hosts, a variety of different lineages evolved which specialized in either bats or birds. The common (C. lectularius) and tropical bed bug (C. hemipterus) split 40 million years before Homo evolution. Humans became hosts to bed bugs through host specialist extension (rather than switching) on three separate occasions.
Historical reports
Bed bugs were first mentioned in ancient Greece as early as 400 BC, and later by Aristotle. Pliny's Natural History, first published circa AD 77 in Rome, claimed bed bugs had medicinal value in treating ailments such as snake bites and ear infections. Belief in the medicinal use of bed bugs persisted until at least the 18th century, when Guettard recommended their use in the treatment of hysteria.
Bed bugs were also mentioned in Germany in the 11th century, in France in the 13th century, and in England in 1583, though they remained rare in England until 1670. Some in the 18th century believed bed bugs had been brought to London with supplies of wood to rebuild the city after the Great Fire of London (1666). Giovanni Antonio Scopoli noted their presence in Carniola (roughly equivalent to present-day Slovenia) in the 18th century.
Traditional methods of repelling or killing bed bugs include the use of plants, fungi, and insects (or their extracts), such as black pepper; black cohosh (Actaea racemosa); Pseudarthria hookeri; Laggera alata (Chinese yángmáo cǎo | 羊毛草); Eucalyptus saligna oil; henna (Lawsonia inermis or camphire); "infused oil of Melolontha vulgaris" (presumably cockchafer); fly agaric (Amanita muscaria); tobacco; "heated oil of Terebinthina" (i.e. true turpentine); wild mint (Mentha arvensis); narrow-leaved pepperwort (Lepidium ruderale); Myrica spp. (e.g. bayberry); Robert geranium (Geranium robertianum); bugbane (Cimicifuga spp.); "herb and seeds of Cannabis"; "opulus" berries (possibly maple or European cranberrybush); masked hunter bugs (Reduvius personatus), "and many others".
In the mid-19th century, smoke from peat fires was recommended as an indoor domestic fumigant against bed bugs.
Dusts have been used to ward off insects from grain storage for centuries, including plant ash, lime, dolomite, certain types of soil, and diatomaceous earth or Kieselguhr. Of these, diatomaceous earth in particular has seen a revival as a non-toxic (when in amorphous form) residual pesticide for bed bug abatement. While diatomaceous earth often performs poorly, silica gel may be effective.
Basket-work panels were put around beds and shaken out in the morning in the UK and in France in the 19th century. Scattering leaves of plants with microscopic hooked hairs around a bed at night, then sweeping them up in the morning and burning them, was a technique reportedly used in Southern Rhodesia and in the Balkans.
Bean leaves have been used historically to trap bedbugs in houses in Eastern Europe. The trichomes on the bean leaves capture the insects by impaling the feet (tarsi) of the insects. The leaves are then destroyed.
20th century
Until the mid-20th century, bed bugs were very common. According to a report by the UK Ministry of Health, in 1933, all the houses in many areas had some degree of bed bug infestation. The increase in bed bug populations in the early 20th century has been attributed to the advent of electric heating, which allowed bed bugs to thrive year-round instead of only in warm weather.
Bed bugs were a serious problem at US military bases during World War II. Initially, the problem was solved by fumigation, using Zyklon Discoids that released hydrogen cyanide gas, a rather dangerous procedure. Later, DDT was used to good effect, though bedbugs have since become largely resistant to it.
The decline of bed bug populations in the 20th century is often credited to potent pesticides that had not previously been widely available. Other contributing factors that are less frequently mentioned in news reports are increased public awareness and slum clearance programs that combined pesticide use with steam disinfection, relocation of slum dwellers to new housing, and in some cases also follow-up inspections for several months after relocated tenants moved into their new housing.
21st century
In 2010, bed bug infestations were reported in New York houses, retail stores, cinemas, offices and schools, especially in Brooklyn and Queens. In early 2023, Orkin reported that Chicago, New York, Philadelphia, Cleveland and Los Angeles were the top five cities in the United States with most bed bug infestations.
In France, these insects re-emerged, despite having disappeared from daily life in the 1950s, due to nomadic lifestyles, consumption of second-hand purchases, and bugs' resistance to insecticides, in addition to increased traveling and tourism following the COVID-19 lockdowns. Between 2017 and 2022, 11% of French households were infested by bed bugs, according to a report from the National Agency for Food, Environmental and Occupational Health Safety (ANSES). In the middle of 2023, reports emerged of a bed bug infestation spread in the capital city of Paris, when it was first seen in cinemas, then it expanded to homes, trains, schools and even hospitals. Treatment of this outbreak has cost France an estimated €230m annually. In the meantime, the United Kingdom witnessed a 65% increase in year-on-year infestations across the country, according to Rentokil. In November 2023, it was reported that South Korea was experiencing a bed bug infestation.
Society and culture
Legal action
Bed bugs are an increasing cause for litigation. Courts have, in some cases, exacted large punitive damage judgments on some hotels. Many of New York City's Upper East Side homeowners have been afflicted, but they tend to remain publicly silent in order not to ruin their property values and be seen as suffering a blight typically associated with "lower social class." Local Law 69 in New York City requires owners of buildings with three or more units to provide their tenants and potential tenants with reports of bedbug history in each unit. They must also prominently post these listings and reports in their building.
Idiom
"Good night, sleep tight, don't let the bed bugs bite," is a traditional saying.
| Biology and health sciences | Hemiptera (true bugs) | null |
952702 | https://en.wikipedia.org/wiki/Anthurium | Anthurium | Anthurium (; Schott, 1829) is a genus of about 1,000 species of flowering plants, the largest genus of the arum family, Araceae. General common names include anthurium, tailflower, flamingo flower, pigtail plant, and laceleaf.
The genus is native to the Americas, where it is distributed from northern Mexico to northern Argentina and parts of the Caribbean.
Description and biology
Anthurium is a genus of herbs often growing as epiphytes on other plants. Some are terrestrial. The leaves are often clustered and are variable in shape. The inflorescence bears small flowers which are perfect, containing male and female structures. The flowers are contained in close together spirals on the spadix. The spadix is often elongated into a spike shape, but it can be globe-shaped or club-shaped. Beneath the spadix is the spathe, a type of bract. This is variable in shape, as well, but it is lance-shaped in many species. It may extend out flat or in a curve. Sometimes it covers the spadix like a hood. The fruits develop from the flowers on the spadix. They are juicy berries varying in color, usually containing two seeds.
The spadix and spathe are a main focus of Anthurium breeders, who develop cultivars in bright colors and unique shapes. Anthurium scherzerianum and A. andraeanum, two of the most common taxa in cultivation, are the only species that grow bright red spathes. They have also been bred to produce spathes in many other colors and patterns.
Anthurium plants are poisonous due to calcium oxalate crystals. The sap is irritating to the skin and eyes.
Cultivation
Like other aroids, many species of Anthurium plant can be grown as houseplants, or outdoors in mild climates in shady spots, including Anthurium crystallinum and Anthurium clarinervium with its large, velvety, dark green leaves and silvery white venation. Many hybrids are derived from Anthurium andraeanum or Anthurium scherzerianum because of their colorful spathes. They thrive in moist soils with high organic matter. In milder climates the plants can be grown in pots of soil. Indoors plants thrive at temperatures of and at lower light than other house plants. Wiping the leaves off with water will remove any dust and insects. Plants in pots with good root systems will benefit from a weak fertilizer solution every other week. In the case of vining or climbing Anthuriums, the plants benefit from being provided with a totem to climb.
Propagation
Anthurium can be propagated by seed or vegetatively by cuttings. In the commercial Anthurium trade, most propagation is via tissue culture.
Species
For a full list, see the List of Anthurium species.
In 1860 there were 183 species known to science, and Heinrich Wilhelm Schott defined them in 28 sections in the book Prodromus Systematis Aroidearum. In 1905 the genus was revised with a description of 18 sections. In 1983 the genus was divided into the following sections:
Belolonchium
Calomystrium
Cardiolonchium
Chamaerepium
Cordatopunctatum
Dactylophyllium
Decurrentia
Digitinervium
Gymnopodium
Leptanthurium
Pachyneurium
Polyphyllium
Polyneurium
Porphyrochitonium
Schizoplacium
Semaeophyllium
Tetraspermium
Urospadix
Xialophyllium
Gallery
Toxicity
All plants within the Anthurium genus are toxic to cats, dogs, and even horses. Each part of the plant, including the root, stems, leaves, flowers, and seeds, poses a risk of toxicity. The plant contains insoluble calcium oxalate crystals, which can cause oral irritation, pain, swelling, excessive drooling, vomiting, and difficulty swallowing. Keeping these plants away from your pets (and equines) is the best way to prevent a medical emergency.
| Biology and health sciences | Monocots | null |
953081 | https://en.wikipedia.org/wiki/Sericulture | Sericulture | Sericulture, or silk farming, is the cultivation of silkworms to produce silk. Although there are several commercial species of silkworms, the caterpillar of the domestic silkmoth is the most widely used and intensively studied silkworm. This species of silkmoth is no longer found in the wild as they have been modified through selective breeding, rendering most flightless and without defense against predators. Silk is believed to have first been produced in China as early as the Neolithic period. Sericulture has become an important cottage industry in countries such as Brazil, China, France, India, Italy, Japan, Korea, Russia, and Thailand. Today, China and India are the two main producers, with more than 60% of the world's annual production.
History
According to Confucian text, the discovery of silk production dates to about 2700 BCE, although archaeological records point to silk cultivation as early as the Yangshao period (5000–3000 BCE). In 1977, a piece of ceramic created 5400–5500 years ago and designed to look like a silkworm was discovered in Nancun, Hebei, providing the earliest known evidence of sericulture. Also, by careful analysis of archaeological silk fibre found on Indus Civilization sites dating back to 2450–2000 BCE, it is believed that silk was being used over a wide region of South Asia. By about the first half of the 1st century CE, it had reached ancient Khotan, by a series of interactions along the Silk Road. By 140 CE, the practice had been established in India. In the 6th century CE, the smuggling of silkworm eggs into the Byzantine Empire led to its establishment in the Mediterranean, remaining a monopoly in the Byzantine Empire for centuries (Byzantine silk). In 1147, during the Second Crusade, Roger II of Sicily (1095–1154) attacked Corinth and Thebes, two important centres of Byzantine silk production, capturing the weavers and their equipment and establishing his own silkworks in Palermo and Calabria, eventually spreading the industry to Western Europe.
Production
The silkworms are fed with mulberry leaves, and after the fourth moult, they climb a twig placed near them and spin their silken cocoons. The silk is a continuous filament comprising fibroin protein, secreted from two salivary glands in the head of each worm, and a gum called sericin, which cements the filaments.
The sericin is removed by placing the cocoons in hot water, which frees the silk filaments and readies them for reeling. This is known as the degumming process. The immersion in hot water also kills the silkmoth pupa.
Single filaments are combined to form thread, in a process called "throwing", which is drawn under tension through several guides and wound onto reels. This process of throwing produces various yarns depending on the amount and direction of the twisting. The threads may be plied to form yarn (short staple lengths are spun; see silk noil). After drying, the raw silk is packed according to quality.
Sustainable silk
Peace silk
The most popular substitute for traditional silk is peace silk, also known as ahimsa silk. The primary factor that makes this form of silk more ethical is that moths are permitted to emerge from their cocoons and fly away prior to boiling. It denotes that no pupa is ever cooked alive during manufacture. However, domesticated silkworms used to make silk have undergone thousands of years of selective breeding and are not "manufactured" to emerge from their cocoons. They are unable to defend themselves against predators since they cannot fly or see clearly. They typically die soon after emerging from their cocoons as a result.
Wild silk
The cocoons of Tussar silkworms, which are found in open woodlands, are used to produce wild silk, also known as Tussar silk. Compared to conventional silk, their cocoons are typically picked after the moths have emerged, making it a more ethical option. Because wild silkworms consume a variety of plants, their fabric is less uniform but more robust. The fabric is made with fewer chemicals as well. The pupae are still inside the cocoons when they are harvested by certain enterprises that employ "wild silk", though.
Stages of production
The stages of production are as follows:
The female silkmoth lays 300 to 500 eggs.
The silkmoth eggs hatch to form larvae or caterpillars, known as silkworms.
The larvae feed on mulberry leaves.
Having grown and moulted several times, the silkworm extrudes a silk fibre and forms a net to hold itself.
It swings itself from side to side in a figure '8', distributing the saliva that will form silk.
The silk solidifies when it contacts the air.
The silkworm spins approximately one mile of filament and completely encloses itself in a cocoon in about two or three days. The amount of usable quality silk in each cocoon is small. As a result, about 2,500 silkworms are required to produce a pound of raw silk.
The intact cocoons are boiled, killing the silkworm pupa.
The silk is obtained by brushing the undamaged cocoon to find the outside end of the filament.
The silk filaments are then wound on a reel. One cocoon contains approximately of silk filament. The silk at this stage is known as raw silk. One thread comprises up to 48 individual silk filaments.
Mahatma Gandhi was critical of silk production based on the Ahimsa philosophy "not to hurt any living thing". He also promoted "Ahimsa silk", made without boiling the pupa to procure the silk and wild silk made from the cocoons of wild and semiwild silkmoths. The Human League also criticised sericulture in their early single "Being Boiled". The organisation PETA has also campaigned against silk.
Pupae as food
The conventional method of silk production results in ~8 kg of wet silkworm pupae and ~2 kg of dry pupae per kilogram of raw silk. This byproduct has historically been consumed by people in silk-producing areas.
Gallery
| Technology | Agriculture_2 | null |
953105 | https://en.wikipedia.org/wiki/Commelina | Commelina | Commelina is a genus of approximately 170 species commonly called dayflowers due to the short lives of their flowers. They are less often known as widow's tears. It is by far the largest genus of its family, Commelinaceae. The Swedish taxonomist Carl Linnaeus of the 18th century named the genus after the two Dutch botanists Jan Commelijn and his nephew Caspar, each representing one of the showy petals of Commelina communis.
The dayflowers are herbs that may be either perennial or annual. They are characterised by their zygomorphic flowers and by the involucral bracts called spathes that surround the flower stalks. These spathes are often filled with a mucilaginous liquid. Each spathe houses either one or two scorpioid cymes, with the upper cyme being either vestigial or bearing from one to several typically male flowers, and the lower cyme bearing several flowers. All members of the genus have alternate leaves.
The Asiatic dayflower (Commelina communis) is probably the best known species in the West. It is a common weed in parts of Europe and throughout eastern North America. Several species, such as Commelina benghalensis, are eaten as a leaf vegetable in Southeast Asia and Africa.
Description
Plants in the genus are perennial or annual herbs with roots that are usually fibrous or rarely tuberous or rhizomatous. The leaves are distichous (i.e. 2-ranked) or spirally arranged with blades that either lack or have a petiole. The ptyxis, or the way the leaf is folded in the bud, is either involute (i.e. having inrolled margins) or supervolute.
The inflorescences are terminal, meaning the stem terminates with an inflorescence, and often leaf opposed, meaning it emerges at the node with a leaf of a new axillary stem. The inflorescence is composed of one or two cincinni, also called scorpioid cymes, which are monochasia (i.e. cymes with a single branched main false axis) in which the lateral branches arise alternately on opposite sides of the false axis. The distal cincinnus may either be vestigial or contain one to several flowers that are typically male. The proximal cyme is always present and is multi-flowered. The cincinni are enclosed in a folded spathe, a modified leaf, which is often filled with a mucilaginous liquid. The spathe may either have completely distinct margins or they may be fused to varying degrees at the basal end.
The flowers are borne on pedicels and are strongly zygomorphic, meaning there is only a single plane of symmetry. Bracteoles occasionally subtend the pedicels, but they are usually absent. The flowers are either bisexual or male. There are three unequal sepals, which may either be free or the two lateral ones may be fused. The petals are free and unequal with the two upper ones being larger and clawed while the lower petal is typically reduced and often differs in colour from the other two. Flower colour is most typically blue, but lilac, lavender, yellow, peach, apricot and white also occur. There are three stamens and two to three staminodes, or infertile stamens, all of which have free, glabrous filaments. The staminodes occur posteriorly and have antherodes with four to six lobes. The stamens are anterior and are longer than the staminodes. The medial stamen differs in size and form from the lateral two, and when a central staminode is present it also differs from the other staminodes. The ovaries are bi- or trilocular and one to two ovules is present per locule.
The fruit is a capsule that is typically bi- or trilocular, but in rare cases may be unilocular, and it is bi- or trivalved. The locules may contain one or two seeds, or no seed at all. The seeds are uniseriate (i.e. arranged in a single row), have a linear hilum and a lateral embryotega.
| Biology and health sciences | Commelinales | Plants |
27850086 | https://en.wikipedia.org/wiki/Small%20modular%20reactor | Small modular reactor | The small modular reactor (SMR) is a class of small nuclear fission reactor, designed to be built in a factory, shipped to operational sites for installation and then used to power buildings or other commercial operations. The term SMR refers to the size, capacity and modular construction. Reactor type and the nuclear processes may vary. Of the many SMR designs, the pressurized water reactor (PWR) is the most common. However, recently proposed SMR designs include: generation IV, thermal-neutron reactors, fast-neutron reactors, molten salt, and gas-cooled reactor models.
Commercial SMRs have been designed to deliver an electrical power output as low as 5 MWe (electric) and up to 300 MWe per module. SMRs may also be designed purely for desalinization or facility heating rather than electricity. These SMRs are measured in megawatts thermal MWt. Many SMR designs rely on a modular system, allowing customers to simply add modules to achieve a desired electrical output.
Small reactors were first designed mostly for military purposes in the 1950s to power ballistic missile submarines and ships (aircraft carriers and ice breakers) with nuclear propulsion. There has been strong interest from technology corporations in using SMRs to power data centers.
The electrical output for modern naval reactors are generally limited to less than 165 MWe and dedicated to powering turboshaft props rather than delivering electricity. In addition, there are many more safety controls absent from naval reactors due to the space limitations these reactors were designed for.
Modular reactors are expected to reduce on-site construction and increase containment efficiency. These reactors are also expected to enhance safety by using passive safety features that do not require human intervention, although this is not specific to SMRs but rather a characteristic of most modern reactor designs.
SMRs are also claimed to have lower power plant staffing costs, as their operation is fairly simple, and are claimed to have the ability to bypass financial and safety barriers that inhibit the construction of conventional reactors.
Working with Oregon State University (OSU), NuScale Power developed the first Nuclear Regulatory Commission approved model for the US market in 2022.
Operational SMRs
As of 2024, only China and Russia have successfully built operational SMRs. There are more than 80 modular reactor designs under development in 19 countries. Russia has been operating a floating nuclear power plant Akademik Lomonosov, in Russia's Far East (Pevek) commercially, since 2020. The floating plant is the first of its kind in the world. China's pebble-bed modular high-temperature gas-cooled reactor HTR-PM was connected to the grid in 2021.
Background
Hope of enhanced safety and reduced costs
Economic factors of scale mean that nuclear reactors tend to be large, to such an extent that size itself becomes a limiting factor. The 1986 Chernobyl disaster and the 2011 Fukushima nuclear disaster caused a major set-back for the nuclear industry, with worldwide suspension of development, cutting down of funding, and closure of reactor plants.
In response, researchers at Oregon State University developed the first commercial SMR prototypes in the late 1990s and early 2000s. A radically different reactor than the one used by the military, OSU's SMR design decreased fabrication time, improved operational safety, and reduced the cost of operation. The goal, of course, was to make it easier for commercial and public entities to afford a traditionally cost-prohibitive form of energy. Credited as the inventor of the commercial SMR, OSU researchers believed the smaller form factor and modular design would allow manufacturers to swap economies-of-unit-scale for economies-of-unit-mass-production - lowering production costs and improving manufacturing efficiency. NuScale Power partnered with OSU to become the first to apply this manufacturing strategy starting in 2006
Proponents claim that SMRs would be less expensive due to the application of standardized modules that could be industrially produced off-site in a dedicated factory. SMRs do, however, also have economic disadvantages. Several studies suggest that the overall costs of SMRs are comparable with those of conventional large reactors. Moreover, extremely limited information about SMR modules transportation has been published. Critics say that modular building will only be cost-effective for a high number of the same SMR type, given the still remaining high costs for each SMR. A high market share is thus needed to obtain sufficient orders.
Contribution to the net zero emissions pathways
In February 2024, the European Commission recognized SMR technology as an important contributor to decarbonization as part of EU Green Deal.
In its pathway to reach global net zero emissions by 2050, the International Energy Agency (IEA) considers that worldwide nuclear power should be multiplied by two between 2020 and 2050. Antonio Vaya Soler, an expert from the Nuclear Energy Agency (NEA), agrees that although renewable energy is essential to fight global warming, it will not be sufficient to achieve net zero emissions and nuclear energy capacity should be at least doubled.
To produce the same electrical power as the nuclear power reactors in the world today, BASE, the German Federal Office for the Safety of Nuclear Waste Management, warns that it would be necessary to build several thousands to tens of thousands of SMRs.
Several fleets of SMRs of exactly the same type, industrially manufactured in large numbers, should be rapidly deployed worldwide to significantly reduce emissions of . The Nuclear Energy Agency (NEA) launched at COP 28 an initiative Accelerating SMRs for Net Zero to foster collaboration between research organizations, nuclear industry, safety authorities, and governments, in order to reduce carbon emissions to net zero before 2050 to limit global surface temperature increase.
Future challenges
Proponents say that nuclear energy with proven technology can be safer; the nuclear industry contends that smaller size will make SMRs even safer than larger conventional plants. This is because the main problem associated with nuclear meltdowns is the decay heat that is present after reactor shutdown, which would be much lower for SMRs because of their lower power output. Critics say that many more small nuclear reactors pose a higher risk, requiring more transportation of nuclear fuel and also increasing the production of radioactive waste. SMRs require new designs with new technology, the safety of which has yet to be proven.
Until 2020, no truly modular SMRs had been commissioned for commercial use. In May 2020, the first prototype of a floating nuclear power plant with two 30 MWe reactors – the type KLT-40 – started operation in Pevek, Russia. This concept is based on the design of nuclear icebreakers. The operation of the first commercial land-based, 125 MWe demonstration reactor ACP100 (Linglong One) is due to start in China by the end of 2026.
Designs
SMRs are envisioned in multiple designs. Some are simplified versions of current reactors, others involve entirely new technologies. All proposed SMRs use nuclear fission with designs including thermal-neutron reactors and fast-neutron reactors.
Thermal-neutron reactors
Thermal-neutron reactors rely on a moderator (water, graphite, beryllium...) to slow neutrons and generally use as fissile material. Most conventional operating reactors are of this type.
Fast reactors
Fast reactors don't use moderators. Instead they rely on the fuel to absorb fast neutrons. This usually means changing the fuel arrangement within the core, or using different fuels. E.g., is more likely to absorb a fast neutron than .
Fast reactors can be breeder reactors. These reactors release enough neutrons to transmute non-fissionable elements into fissionable ones. A common use for a breeder reactor is to surround the core by a "blanket" of , the most easily available isotope. Once the undergoes a neutron absorption reaction, it becomes , which can be removed from the reactor during refueling, and subsequently reprocessed and used as fuel.
Technologies
Coolant
Conventional light-water reactors typically use water as a coolant and neutron moderator. SMRs may use water, liquid metal, gas and molten salt as coolants. Coolant type is determined based on the reactor type, reactor design, and the chosen application. Large-rated reactors primarily use light water as coolant, allowing for this cooling method to be easily applied to SMRs. Helium is often elected as a gas coolant for SMRs because it yields a high plant thermal efficiency and supplies a sufficient amount of reactor heat. Sodium, lead, and lead-bismuth eutectic (LBE) are liquid metal coolants studied for 4th generation SMRs. There was a large focus on sodium during early work on large-rated reactors which has since carried over to SMRs to be a prominent choice as a liquid metal coolant. SMRs have lower cooling water requirements, which expands the number of sites where a SMR could be built, including remote areas typically incorporating mining and desalination.
Thermal/electrical generation
Some gas-cooled reactor designs could drive a gas turbine, rather than boiling water, such that thermal energy can be used directly. Heat could also be used in hydrogen production and other industrial operations, such as desalination and the production of petroleum derivative (extracting oil from oil sands, making synthetic oil from coal, etc.).
Load following
SMR designs are generally expected to provide base load electrical power; some proposed designs are aimed to adjust their power output based on electricity demand.
Another approach, especially for SMRs designed to provide high temperature heat, is to adopt cogeneration, maintaining consistent heat output, while diverting otherwise unneeded heat to an auxiliary use. District heating, desalination and hydrogen production have been proposed as cogeneration options.
Overnight desalination requires sufficient freshwater storage capacity to deliver water at times other than when it is produced. Reverse osmosis membrane and thermal evaporators are the two main techniques for seawater desalination. The membrane desalination process uses only electricity to power water pumps and is the most employed of the two methods. In the thermal process, the feed water stream is evaporated in different stages with continuous decreases in pressure between the stages. The thermal process directly uses thermal energy and avoids the conversion of thermal power into electricity. Thermal desalination is further divided into two main technologies: the multi-stage flash distillation (MSF) and the Multi-Effect Desalination (MED).
Nuclear safety
A report by the German Federal Office for the Safety of Nuclear Waste Management (BASE) considering 136 different historical and current reactors and SMR concepts stated: "Overall, SMRs could potentially achieve safety advantages compared to power plants with a larger power output, as they have a lower radioactive inventory per reactor and aim for a higher safety level especially through simplifications and an increased use of passive systems. In contrast, however, various SMR concepts also favour reduced regulatory requirements, for example, with regard to the required degree of redundancy or diversity in safety systems. Some developers even demand that current requirements be waived, for example in the area of internal accident management or with reduced planning zones, or even a complete waiver of external emergency protection planning. Since the safety of a reactor plant depends on all of these factors, based on the current state of knowledge it is not possible to state, that a higher safety level is achieved by SMR concepts in principle."
Negative temperature coefficients in the moderators and the fuels keep the fission reactions under control, causing the reaction to slow as temperature increases. After the shutdown of a nuclear reactor, the reactor needs to be cooled continuously in order to dissipate decay heat. A loss of emergency cooling such as in the Fukushima nuclear accident and the Three Mile Island accident can result in a nuclear meltdown when the temperature in the reactor becomes too high. Since the initial decay heat is a fraction of the reactor operating power, the lower operating power of SMRs makes them much safer since less heat needs to be dissipated.
Some SMR designs proposes cooling systems only based on thermoconvection – natural circulation – to eliminate cooling pumps that could break down. Convection can keep removing decay heat after reactor shutdown. However, some SMRs may need an active cooling system to back up the passive system, increasing cost.
Some SMR designs feature an integral design of which the primary reactor core, steam generator and the pressurizer are integrated within the sealed reactor vessel. This integrated design allows for the reduction of a possible accident as contamination leaks could be contained. In comparison to larger reactors having numerous components outside the reactor vessel, this feature increases the safety by decreasing the risks of an uncontained accident. Some SMR designs also envisage to install the reactor and the spent-fuel storage pools underground.
Radioactive waste
The backend of the nuclear fuel cycle of SMR is a complex and challenging issue remaining disputed. The quantity and the radiotoxicity of the radioactive waste produced by SMR mainly depend on their design and the related fuel cycle. As the notion of SMR encompasses a broad spectrum of nuclear reactor types, no simple answer can be easily given to the question. SMR may comprise small light water reactors of third generation as well as small fast neutron reactors of fourth generation.
Often, the startup companies developing unconventional SMR prototypes advocate waste reduction as an advantage of the proposed solution and even sometimes claim that their technology could eliminate the need for a deep geological repository to dispose of high-level and long-lived radioactive waste. This is especially the case for companies studying fast neutron reactors of 4th generation (molten salts reactors, metal-cooled reactors (sodium-cooled fast reactor, or lead-cooled fast reactor).
Fast breeder reactors "burn" (0.7% of natural uranium), but also convert fertile materials such as (99.3% of natural uranium) into fissile that can be used as nuclear fuel.
The traveling wave reactor proposed by TerraPower is aimed to immediately "burn" the fuel that it breeds without requiring its removal from the reactor core and its further reprocessing.
The design of some SMR reactors is based on the thorium fuel cycle, which is considered by their promotors as a way to reduce the long-term waste radiotoxicity compared to the uranium cycle. However, using the thorium cycle also presents big operational challenges because of the production and the use of and long-lived fertile , both radioisotopes emitting strong gamma rays. So, the presence of these radionuclides seriously complicates the radiation shielding of the fresh nuclear fuel and the long-term storage and disposal of their spent nuclear fuel.
A study of 2022 made by Krall, Macfarlane and Ewing is more critical and reports that some types of SMR could produce more waste per unit of output power than conventional reactors, in some cases more than 5× the amount of spent fuel per kilowatt, and as much as 35× for other waste produced by neutron activation, such as activated steel and graphite.
These authors have identified the neutron leakage as the first issue for SMRs because they have a higher surface area with respect to their core volume. They have calculated that the neutron leakage rates are much higher for SMRs, because in smaller reactor cores, emitted neutrons have fewer chances to interact with the fissile atoms present in the fuel and to produce nuclear fission. Instead, neutrons exit the reactor core without interacting with the nuclear fuel, and they are absorbed outside the core by the materials used for the neutron reflectors and the shielding (thermal and gamma shields), turning them as radioactive waste (activated steel and graphite).
Reactor designs using liquid metal coolants (molten sodium, lead, lead-bismuth eutectic, LBE) also become radioactive and contains activated impurities.
Another issue pinpointed by Krall et al. (2022) related to the higher neutron leakage in SMR is that a lower fraction of their nuclear fuel is consumed, leading to a lower burnup and to more fissile materials left over in their spent fuel, therefore increasing the waste volume. To sustain the nuclear chain reactions in the core of a smaller reactor, an alternative is to use nuclear fuel more enriched in . This could increase the risks of nuclear proliferation and could require more stringent safeguard measures to prevent it (see also IAEA safeguards).
If higher concentrations of fissile materials subsist in the spent fuel, the critical mass needed to sustain a nuclear chain reaction is also lower. As a direct consequence, the number of spent fuels present in a waste canister will also be lower and a larger number of canisters and overpacks will be necessary to avoid criticality accidents and to guarantee nuclear criticality safety in a deep geological repository. This also contributes to increase the total waste volume and the number of disposal galleries in a geological repository.
Given the potential technical and economical importance of SMRs to supply zero-carbon electrical energy needed to fight climate change and the long-term and social relevance of the study to adequately manage and dispose of radioactive waste without imposing a negative burden onto the future generations, the publication of Krall et al. (2022) in the prestigious PNAS journal has attracted many reactions ranging from criticisms on the quality of their data and hypotheses to international debates on radioactive waste produced by SMRs and their decommissioning.
In an interview with François Diaz-Maurin, the associate editor of the Bulletin of the Atomic Scientists, Lindsay Krall, the lead author of the study and a former MacArthur postdoctoral fellow at Stanford's Center for International Security and Cooperation (CISAC) answered to questions and criticisms, amongst others, those raised by the NuScale reactor company. One of the main concerns Krall expressed in this interview is that:
"There's definitely a disconnect between the people working on the back end of the fuel cycle—especially with geologic repository development—and those actually designing reactors. And, there is not a lot of motivation for these reactor designers to think about the geologic disposal aspects because the NRC's new reactor design certification application does not have a chapter on geologic disposal..."
The critical study of Krall et al. (2022) has the merit to have raised relevant questions that cannot be ignored by reactor designers, or decision-makers, and to have triggered open and fresh discussions on important outcomes for SMRs and radwaste management in general. Amongst the various types of SMR projects initiated today by many start-up companies, only those correctly addressing these questions and really contributing to minimize the radioactive waste they produce have a chance to be supported by the public and governmental organisations (nuclear safety authorities and radioactive waste management organisations) and their research to be funded by long-term national policies.
The high diversity of SMR reactors and their respective fuel cycles may also require more diverse waste management strategy to recycle, or to safely dispose, their nuclear waste. A larger number of spent fuel types will be more difficult to manage than only one type as it is presently the case with light water reactors only.
As previously stressed by Krall and Macfarlane (2018), some types of SMR spent fuels, or coolants, (highly reactive and corrosive uranium fluoride () from molten salt reactors, or pyrophoric sodium from liquid metal cooled fast breeders) cannot be directly disposed of in a deep geologic repository because of their chemical reactivity in the underground environment (deep clay formations, crystalline rocks, or rock salt). To avoid to exacerbate spent fuel storage and disposal issues it will be mandatory to reprocess and to condition them in an appropriate and safe way before final geological disposal.
A study made by Keto et al. (2022) at the VTT Technical Research Centre of Finland also addressed the management of spent nuclear fuel (SNF) and low- and intermediate-level waste (LILW) from a possible future deployment of SMRs in Finland. It also indicates that larger masses (per GWe-year) of SNF and other HLW and larger volumes (per GWe-year) of LLW would be produced by a light water SMR compared to a large NPP.
A report by the German Federal Office for the Safety of Nuclear Waste Management (BASE) found that extensive interim storage and fuel transports are still required for SMRs. A deep geological repository is still unavoidable in any case because of the presence of highly mobile long-lived fission products that, due to their too low neutron cross section, cannot be efficiently transmuted, as it is the case with dose-dominating radionuclides such as , and (soluble anions that are not sorbed onto the negatively charged minerals and are not retarded in geological media).
Nuclear proliferation
Nuclear proliferation, or the use of nuclear materials to create weapons, is a concern for small modular reactors. As SMRs have lower generation capacity and are physically smaller, they are intended to be deployed in many more locations than conventional plants. SMRs are expected to substantially reduce staffing levels. The combination creates physical protection and security concerns.
SMRs can be designed to use unconventional fuels allowing for higher burnup and longer fuel cycles. Longer refueling intervals could contribute to decrease the proliferation risks. Once the fuel has been irradiated, the mixture of fission products and fissile materials is highly radioactive and requires special handling, preventing casual theft.
Contrasting to conventional large reactors, SMRs can be adapted to be installed in a sealed underground chamber; therefore, "reducing the vulnerability of the reactor to a terrorist attack or a natural disaster". New SMR designs enhance the proliferation resistance, such as those from the reactor design company Gen4. These models of SMR offer a solution capable of operating sealed underground for the life of the reactor following installation.
Some SMR designs are designed for one-time fueling. This improves proliferation resistance by eliminating on-site nuclear fuel handling and means that the fuel can be sealed within the reactor. However, this design requires large amounts of fuel, which could make it a more attractive target. A 200 MWe 30-year core life light water SMR could contain about 2.5 tonnes of plutonium at end of life.
Furthermore, many SMRs offer the ability to go periods of greater than 10 years without requiring any form of refueling therefore improving the proliferation resistance as compared to conventional large reactors of which entail refueling every 18–24 months.
Light-water reactors designed to run on thorium offer increased proliferation resistance compared to the conventional uranium cycle, though molten salt reactors have a substantial risk.
SMRs are transported from the factories without fuel, as they are fueled on the ultimate site, except some microreactors. This implies an independent transport of the fuel to the site and therefore increases the risk of nuclear proliferation.
Licensing process
Licensing is an essential process required to guarantee the safety, security and safeguards of a new nuclear installation. Only NuScale Power's VOYGR SMR is fully licensed for use in the United States. However, not all countries follow the NRC or IAEA licensing standards. In the United States and IAEA adhering countries, the licensing is based on a rigorous, independent analysis and reviewing work of all structures, systems and components critical for the nuclear safety under normal and accidental conditions on the whole service life of the installation including the long-term management of radioactive waste. Licensing is based on the examination and scrutiny of the risk assessment studies and safety files elaborated by the fabricant and the exploitant of the SMR in the frame of the safety case they have to submit to the safety authority (regulatory body) when applying for a licence to construct and safely exploit the installation. For NRC and IAEA licensing, the safety and feasibility cases of nuclear installations have to take into account all processes and elements important for the operational safety, its security (access protection), the nuclear safeguard (risk of proliferation), the proper conditioning of radioactive waste under a stable physico-chemical form, and the long-term safety related to the final disposal of the different types of radwaste produced, including all the waste produced during dismantling operations after decommissioning of the installation. A particularly important point of attention for the backend of the nuclear fuel cycle is to avoid to producing poorly conditioned waste, or waste types without sustainable final destination or susceptible to generating unexpected reprocessing and disposal costs.
The most common licensing process, applied by existing commercial reactors, is for the operation of light water reactors (PWR and BWR). Early designs for large-scale reactors date back to the 1960s and 1970s during the construction of the nuclear reactor fleet currently in service. Some adaptations of the original licensing process by the US's Nuclear Regulatory Commission (NRC) have been repurposed to better correspond to the specific characteristics and needs of the deployment of SMR units. In particular, the US Nuclear Regulatory Commission process for licensing has focused mainly on conventional reactors. Design and safety specifications, human and organizational factors (including staffing requirements) have been developed for reactors with electrical output of more than 700 MWe.
To ensure adequate guidelines for the nuclear safety, while helping the licensing process, the IAEA has encouraged the creation of a central licensing system for SMRs. A workshop in October 2009 and another in June 2010 considered the topic, followed by an US congressional hearing in May 2010.
The NRC and the United States Department of Energy) are working to define SMR licensing. The challenge of facilitating the development of SMRs is to prevent a weakening of the safety regulations: the risk of lightened regulations adopted more rapidly is to lower the safety characteristics of SMRs. While deploying identical systems built in manufacturing plants with an improved quality control can be considered an advantage, SMRs remain nuclear reactors with a very high energy density and their smaller size is not per se an intrinsic guarantee for a better safety. Any severe accident with external radioactive contamination release could have potential serious consequences not so different from that of a large LWR reactor. It would also probably signify the final rejection of nuclear energy by the public and the end of the nuclear industry. The potential "proliferation" of large SMR fleets and the high diversity of their design also complicate the licensing process. The nuclear safety cannot be sacrificed for industrial or economical interests and the risk of nuclear accident increases with the number of reactors in service, small or large unit.
The U.S. Advanced Reactor Demonstration Program was expected to help license and build two prototype SMRs during the 2020s, with up to $4 billion of government funding.
In July 2024, the ADVANCE Act directed the United States Nuclear Regulatory Commission to develop a process to license and regulate microreactor designs. The Act is intended to expedite the deployment of microreactors, among other nuclear technologies.
Flexibility
Small nuclear reactors, in comparison to conventional nuclear power plants, offer potential advantages related to the flexibility of their modular construction. It would be possible to incrementally connect additional units to the grid in the event electrical load increases. Additionally, this flexibility in a standardized SMRs design revolving around modularity could allow for a faster production at a decreasing cost following the completion of the first reactor on site.
The hypothesised flexibility and modularity of SMR is intended to allow additional power generation capability to be installed at existing power plants. A site could host several SMRs, one going off-line for refueling while the other reactors stay online as it is presently already the case for conventional larger reactors.
When electrical energy is not needed, some SMR designs foresee the direct use of thermal energy, minimizing so the energy loss. This includes "desalination, industrial processes, hydrogen production, shale oil recovery, and district heating", uses for which the present conventional larger reactors are not designed.
Economics
A key driver of interest in SMRs is the claimed economies of scale in production, due to volume manufacture in an offsite factory. Some studies instead find the capital cost of SMRs to be equivalent to larger reactors. Substantial capital is needed to construct the factory – ameliorating that cost requires significant volume, estimated to be 40–70 units.
Another potential advantage is that a future power station using SMRs can begin with a single module and expand by adding modules as demand grows. This reduces startup costs associated with conventional designs. Some SMRs also have a load-following design such that they could produce less electricity when demand is low.
According to a 2014 study of electricity production in decentralized microgrids, the total cost of using SMRs for electricity generation would be significantly lower compared to the total cost of offshore wind power, solar thermal energy, biomass, and solar photovoltaic electricity generation plants.
Construction costs per SMR reactor were claimed in 2016 to be less than that for a conventional nuclear plant, while exploitation costs might be higher for SMRs due to low scale economics and the higher number of reactors. SMR staff operating costs per unit output can be as much as 190% higher than the fixed operating cost of fewer large reactors. Modular building is a very complex process and there is "extremely limited information about SMR modules transportation", according to a 2019 report.
A production cost calculation done by the German Federal Office for the Safety of Nuclear Waste Management (BASE), taking into account economies of scale and learning effects from the nuclear industry, suggests that an average of 3,000 SMR would have to be produced before SMR production would be worthwhile. This is because the construction costs of SMRs are relatively higher than those of large nuclear power plants due to the low electrical output.
In 2017, an Energy Innovation Reform Project (EIRP) study of eight companies looked at reactor designs with capacity between 47.5 MWe and 1,648 MWe. The study reported average capital cost of $3,782/kW, average operating cost total of $21/MWh and levelized cost of electricity (LCOE) of $60/MWh.
In 2020, Energy Impact Center founder Bret Kugelmass claimed that thousands of SMRs could be built in parallel, "thus reducing costs associated with long borrowing times for prolonged construction schedules and reducing risk premiums currently linked to large projects". GE Hitachi Nuclear Energy Executive Vice President Jon Ball agreed, saying the modular elements of SMRs would also help reduce costs associated with extended construction times.
In October 2023, an academic paper published in Energy collated the basic economic data of 19 more developed SMR designs, and modeled their costs in a consistent manner. A Monte Carlo simulation showed that none were profitable or economically competitive. For the closer to market PWR SMRs the median LCOEs ranged from $218/MWh to $614/MWh (in 2020 US dollars), with lower first quartile estimates from $188/MWh to $385/MWh. The three high-temperature gas-cooled reactor designs, which needed more development time, had lower median LCOEs from $116/MWh to $137/MWh.
The first SMR deployment project in the US was the Carbon Free Power Project, which planned to deploy six 77 MWe NuScale reactors, reduced from twelve in earlier plans. Estimated target electricity generation price after subsidies was $89/MWh in 2023, an increase from $58/MWh in 2021. The increased generation cost led to the decision to cancel the project in November 2023. Before its cancellation, the project received a $1.355 billion cost-share award toward construction costs from the US government in 2020 plus an estimated $30/MWh generation subsidy from the 2020 Inflation Reduction Act. Unsubsidized cost estimates at cancellation were a capital cost of $20,139/kW and generating cost of $119/MWe. This raised concerns about the commercial prospects in the U.S. of the other SMR designs.
In 2024, Australian scientific research body CSIRO estimated that electricity produced in Australia by a SMR constructed from 2023 would cost roughly 2.5 times that produced by a traditional large nuclear plant, falling to about 1.6 times by 2030.
List of reactor designs
Numerous reactor designs have been proposed. Notable SMR designs:
Siting/infrastructure
SMRs are expected to require less land, e.g., the 470 MWe 3-loop Rolls-Royce SMR reactor should take , 10% of that needed for a traditional plant. This unit is too large to meet the International Atomic Energy Agency's definition of a SMR being smaller than 300MWe and will require more on-site construction, which calls into question the claimed benefits of SMRs. The firm is targeting a 500-day construction time.
Electricity needs in remote locations are usually small and variable, making them suitable for a smaller plant. The smaller size may also reduce the need to access to a large grid to distribute their output.
Proposed sites
Argentina
In February 2014, the CAREM SMR project started in Argentina with the civil engineering construction of the containment building of a prototype reactor. The CAREM acronym means . The National Atomic Energy Commission (, CNEA), the Argentine government agency in charge of nuclear energy research and development and , the national nuclear energy company, are cooperating to achieve the realization of the project.
CAREM-25 is a prototype of 25 MWe, the first nuclear power plant completely designed and developed in Argentina. The project was suspended several times before being resumed. In October 2022, CNEA expected that the civil construction works would be finished by 2024. If construction continues according to plan, the first criticality of CAREM-25 is foreseen by the end of 2027.
Canada
In 2018, the Canadian province of New Brunswick announced it would invest $10 million for a demonstration project at the Point Lepreau Nuclear Generating Station. It was later announced that SMR proponents Advanced Reactor Concepts and Moltex would open offices there.
One unit is scheduled for construction at Point Lepreau Nuclear Generating Station, Canada, in July 2018. Both Moltex and ARC Nuclear are vying for the contract.
On 1 December 2019, the Premiers of Ontario, New Brunswick and Saskatchewan signed a memorandum of understanding (MoU) "committing to collaborate on the development and deployment of innovative, versatile and scalable nuclear reactors, known as Small Modular Reactors (SMRs)." They were joined by Alberta in August 2020. With continued support from citizens and government officials have led to the execution of a selected SMR at the Canadian Nuclear Laboratory.
In 2021, Ontario Power Generation announced they plan to build a BWRX-300 SMR at their Darlington site to be completed by 2028. A licence for construction still had to be applied for.
On 11 August 2022, Invest Alberta, the Government of Alberta's crown corporation signed a MoU with Terrestrial Energy regarding IMSR in Western Canada through an interprovincial MoU it joined earlier.
China
In July 2019, China National Nuclear Corporation announced it would build an ACP100 SMR on the north-west side of the existing Changjiang Nuclear Power Plant at Changjiang, in the Hainan province by the end of the year. On 7 June 2021, the demonstration project, named the Linglong One, was approved by China's National Development and Reform Commission. In July, China National Nuclear Corporation (CNNC) started construction, and in October 2021, the containment vessel bottom of the first of two units was installed. It is the world's first commercial land-based SMR prototype.
In August 2023, the core module was installed. The core module includes an integrated pressure vessel, steam generator, primary pump receiver. The reactor's planned capacity is 125 MWe.
France
At the beginning of 2023, Électricité de France (EDF) created a new subsidiary to develop and construct a new SMR named Nuward. It was a 340 MWe design with two independent light water reactors of 170 MWe. The twin reactors were sheltered in a single containment building sharing most of their equipment. In August 2023, EDF submitted a safety case for Nuward to the (ASN), the French safety authority.
In July 2024, EDF announced it was discontinuing the existing design process for Nuward, and will work on an SMR design based on existing rather than innovative technologies, following discussions with prospective SMR customers. In January 2025, EDF announced that the new Nuward conceptual design would be completed by mid-2026 to come to market in the 2030s, with an output of about 400 MWe and usable heat output of 100 MWt.
Poland
Polish chemical company Synthos declared plans to deploy a Hitachi BWRX-300 reactor (300 MW) in Poland by 2030. A feasibility study was completed in December 2020 and the licensing process started with the Polish National Atomic Energy Agency.
In February 2022, NuScale Power and the large mining conglomerate KGHM Polska Miedź announced signing of contract to construct a first operational reactor in Poland by 2029.
Romania
On the occasion of 2021 United Nations Climate Change Conference, the state-owned Romanian nuclear energy company Nuclearelectrica and NuScale Power signed an agreement to build a power plant with six small-scale nuclear reactors at the Doicești power station, on the site of a former coal power plant, located near the village of Doicești, Dâmbovița county, 90 km North of Bucharest. The project is estimated to be completed by 2026–2027, which will make the power plant the first of its kind in Europe. The power plant is expected to generate 462 MWe, securing the consumption of about 46.000 households and would help to avoid the release of 4 million tons of per year.
Russia
Russia has started to deploy on its arctic coast small nuclear reactors embarked on board icebreakers. In May 2020, the first prototype of a floating nuclear power plant with two 30 MWe reactors – the type KLT-40 – started operation in Pevek, Russia. This concept is based on the design of nuclear icebreakers.
United Kingdom
In 2016, it was reported that the UK Government was assessing Welsh SMR sites – including the former Trawsfynydd nuclear power station – and on the site of former nuclear or coal-fired power stations in Northern England. Existing nuclear sites including Bradwell, Hartlepool, Heysham, Oldbury, Sizewell, Sellafield, and Wylfa were stated to be possibilities. The target cost for a 470 MWe Rolls-Royce SMR unit is £1.8 billion for the fifth unit built. In 2020, it was reported that Rolls-Royce had plans to construct up to 16 SMRs in the UK. In 2019, the company received £18 million to begin designing the modular system. An additional £210 million was awarded to Rolls-Royce by the British government in 2021, complemented by a £195 million contribution from private firms. In November 2022, Rolls-Royce announced that the sites at Trawsfynydd, Wylfa, Sellafield and Oldbury would be prioritised for assessment as potential locations for multiple SMRs.
The British government launched Great British Nuclear in July 2023 to administer a competition to create SMRs, and will co-fund any viable project.
United States
The US Department of Energy had estimated the first SMR in the United States would be completed by NuScale Power around 2030, but this deal has since fallen through after the customers backed out due to rising costs. The United States has plans for several modular reactors. Dominion Energy Virginia is now accepting proposals. The U.S. has nearly 4 gigawatts in announced SMR projects in addition to almost 3 GW in early development or pre-development stages, according to Utility Dive.
SMRs differ in terms of staffing, safety and deployment time. US government studies to evaluate SMR-associated risks are claimed to have slowed the licensing process. One main concern with SMRs and their large number, needed to reach an economic profitability, is preventing nuclear proliferation.
Standard Power, a provider of infrastructure as a service to advanced data processing companies, has chosen to work with NuScale Power and ENTRA1 Energy to develop SMR-powered facilities in Pennsylvania and Ohio that will together produce nearly two gigawatts of clean, reliable energy.
NuScale Power is working with Wisconsin's Dairyland Power to evaluate VOYGR SMR power plants for potential deployment. The US leader in SMR technology believes its load-following capabilities can be used to support Dairyland's existing renewables portfolio, as well as facilitate growth. Additionally, VOYGR plants are well-suited for replacing Dairyland's retiring coal plant sites, preserving critical jobs and helping communities transition to a decarbonized energy system.
NuScale Power is working with Associated Electric Cooperative Inc. (Associated) in Missouri to evaluate deployment of VOYGR SMR power plants as part of Associated's due diligence to explore reliable, responsible sources of energy.
The Utah Associated Municipal Power Systems (UAMPS) had partnered with Energy Northwest to explore siting a NuScale Power reactor in Idaho, possibly on the Department of Energy's Idaho National Laboratory. Known as the Carbon Free Power Project, the project was canceled in November 2023 for cost reasons. NuScale said in January 2023 the target price for power from the plant was $89 per megawatt hour, up 53% from the previous estimate of $58 per MWh, raising concerns about customers' willingness to pay. Still, increased cost estimates remain well below traditional nuclear power used for commercial facilities and most other less reliable and more environmentally hazardous forms of power production.
The Galena Nuclear Power Plant in Galena, Alaska was a proposed micro nuclear reactor installation. It was a potential deployment for the Toshiba 4S reactor. The project was "effectively stalled". Toshiba never began the expensive process for approval that is required by the U.S. Nuclear Regulatory Commission.
Although the SMR now under consideration has yet to be NRC licensed, the Tennessee Valley Authority was authorized to receive an Early Site Permit (ESP) by the Nuclear Regulatory Commission for siting an SMR at its Clinch River Nuclear Site in Tennessee in December 2019. This ESP is valid for 20 years, and addresses site safety, environmental protection and emergency preparedness. This ESP is applicable for any light-water reactor SMR design under development in the United States.
In October 2024, Google agreed to commission multiple small modular reactors from Kairos Power to power its artificial intelligence processing, with the first to be operational in 2030.
| Technology | Power generation | null |
5275240 | https://en.wikipedia.org/wiki/Photon%20gas | Photon gas | In physics, a photon gas is a gas-like collection of photons, which has many of the same properties of a conventional gas like hydrogen or neon – including pressure, temperature, and entropy. The most common example of a photon gas in equilibrium is the black-body radiation.
Photons are part of a family of particles known as bosons, particles that follow Bose–Einstein statistics and with integer spin. A gas of bosons with only one type of particle is uniquely described by three state functions such as the temperature, volume, and the number of particles. However, for a black body, the energy distribution is established by the interaction of the photons with matter, usually the walls of the container, and the number of photons is not conserved. As a result, the chemical potential of the black-body photon gas is zero at thermodynamic equilibrium. The number of state variables needed to describe a black-body state is thus reduced from three to two (e.g. temperature and volume).
Thermodynamics of a black body photon gas
In a classical ideal gas with massive particles, the energy of the particles is distributed according to a Maxwell–Boltzmann distribution. This distribution is established as the particles collide with each other, exchanging energy (and momentum) in the process. In a photon gas, there will also be an equilibrium distribution, but photons do not collide with each other (except under very extreme conditions, see two-photon physics), so the equilibrium distribution must be established by other means. The most common way that an equilibrium distribution is established is by the interaction of the photons with matter. If the photons are absorbed and emitted by the walls of the system containing the photon gas, and the walls are at a particular temperature, then the equilibrium distribution for the photons will be a black-body distribution at that temperature.
A very important difference between a generic Bose gas (gas of massive bosons) and a photon gas with a black-body distribution is that the number of photons in the photon gas is not conserved. A photon can be created upon thermal excitation of an atom in the wall into an upper electronic state, followed by the emission of a photon when the atom falls back to a lower energetic state. This type of photon generation is called thermal emission. The reverse process can also take place, resulting in a photon being destroyed and removed from the gas. It can be shown that, as a result of such processes there is no constraint on the number of photons in the system, and the chemical potential of the photons must be zero for black-body radiation.
The thermodynamics of a black-body photon gas may be derived using quantum statistical mechanical arguments, with the radiation field being in equilibrium with the atoms in the wall. The derivation yields the spectral energy density u, which is the energy of the radiation field per unit volume per unit frequency interval, given by:
.
where h is the Planck constant, c is the speed of light, ν is the frequency, k is the Boltzmann constant, and T is temperature.
Integrating over frequency and multiplying by the volume, V, gives the internal energy of a black-body photon gas:
.
The derivation also yields the (expected) number of photons N:
,
where is the Riemann zeta function. Note that for a particular temperature, the particle number N varies with the volume in a fixed manner, adjusting itself to have a constant density of photons.
If we note that the equation of state for an ultra-relativistic quantum gas (which inherently describes photons) is given by
,
then we can combine the above formulas to produce an equation of state that looks much like that of an ideal gas:
.
The following table summarizes the thermodynamic state functions for a black-body photon gas. Notice that the pressure can be written in the form , which is independent of volume (b is a constant).
Isothermal transformations
As an example of a thermodynamic process involving a photon gas, consider a cylinder with a movable piston. The interior walls of the cylinder are "black" in order that the temperature of the photons can be maintained at a particular temperature. This means that the space inside the cylinder will contain a blackbody-distributed photon gas. Unlike a massive gas, this gas will exist without the photons being introduced from the outside – the walls will provide the photons for the gas. Suppose the piston is pushed all the way into the cylinder so that there is an extremely small volume. The photon gas inside the volume will press against the piston, moving it outward, and in order for the transformation to be isothermic, a counter force of almost the same value will have to be applied to the piston so that the motion of the piston is very slow. This force will be equal to the pressure times the cross sectional area () of the piston. This process can be continued at a constant temperature until the photon gas is at a volume V0. Integrating the force over the distance () traveled yields the total work done to create this photon gas at this volume
,
where the relationship V = Ax has been used. Defining
.
The pressure is
.
Integrating, the work done is just
.
The amount of heat that must be added in order to create the gas is
.
where H0 is the enthalpy at the end of the transformation. It is seen that the enthalpy is the amount of energy needed to create the photon gas.
Photon gases with tunable chemical potential
In low-dimensional systems, for example in dye-solution filled optical microcavities with a distance between the resonator mirrors in the wavelength range where the situation becomes two-dimensional, also photon gases with tunable chemical potential can be realized. Such a photon gas in many respects behaves like a gas of material particles. One consequence of the tunable chemical potential is that at high phase space densities then Bose-Einstein condensation of photons is observed.
| Physical sciences | Thermodynamics | Physics |
34693203 | https://en.wikipedia.org/wiki/Terrestrial%20crab | Terrestrial crab | A number of lineages of crabs have evolved to live predominantly on land. Examples of terrestrial crabs are found in the families Gecarcinidae and Gecarcinucidae, as well as in selected genera from other families, such as Sesarma, although the term "land crab" is often used to mean solely the family Gecarcinidae.
Terrestriality and migration
No clear distinction is made between "terrestrial", "semiterrestrial", and "aquatic" crabs. Rather, a continuum of terrestriality is displayed among the true crabs, although most land-adapted crabs must still return to water to release their eggs. Some species of terrestrial crabs can be found many kilometres from the sea, but have to complete annual migrations to the sea. For example, following the Indian Ocean monsoon, the Christmas Island red crab (Gecarcoidea natalis) migrates , forming a "living carpet" of crabs. The crabs can travel up to in a day, and up to in total. Only a few land crabs, including certain Geosesarma species, have direct development (the mother carries the eggs until they have become tiny, fully developed crabs), and these do not need access to water to breed. Many crabs belonging to the family Potamidae, which contains mostly freshwater crabs, have developed a semiterrestrial (for instance the genus Nanhaipotamon) to terrestrial life history, and are sometimes independent of fresh water for reproduction (for instance the genus Tiwaripotamon).
Ecology
Terrestrial crabs are often similar to freshwater crabs, since the physiological changes needed for living in fresh water are preadaptations for terrestrial living. On some oceanic islands, terrestrial crabs occupy the top of the energy pyramid.
| Biology and health sciences | Crabs and hermit crabs | Animals |
1450654 | https://en.wikipedia.org/wiki/Lablab | Lablab | Lablab purpureus is a species of bean in the family Fabaceae. It is native to sub-Saharan Africa and India and it is cultivated throughout the tropics for food. English language common names include hyacinth bean, lablab-bean bonavist bean/pea, dolichos bean, seim or sem bean, lablab bean, Egyptian kidney bean, Indian bean, bataw and Australian pea. Lablab is a monotypic genus.
Taxonomy
The name lablab which is also capitalized as its genus name is given by Robert Sweet from the previous name of Dolichos lablab by Carl Linnaeus, its epithet comes from .
Subspecific classification
According to the British biologist and taxonomist Bernard Verdcourt,
there are two cultivated subspecies of Lablab purpureus (L.) Sweet:
Lablab purpureus subsp. bengalensis (Jacq.) Verdc. (Syn.: Dolichos bengalensis Jacq., Dolichos lablab subsp. bengalensis (Jacq.) Rivals, Lablab niger subsp. bengalensis (Jacq.) Cuf.)
Lablab purpureus subsp. purpureus
in addition to one wild subspecies:
Lablab purpureus subsp. uncinatus
of which a special variant with lobed leaflets exists only in Namibia:
Lablab purpureus var. rhomboïdeus (Schinz).
Description
The plant is variable due to extensive breeding in cultivation, but in general, they are annual or short-lived perennial vines. The wild species is perennial. The thick stems can reach in length. The leaves are made up of three pointed leaflets, each up to long. They may be hairy on the undersides. The inflorescence is made up of racemes of many flowers. Some cultivars have white flowers, and others may have purplish or blue.
The fruit is a legume pod variable in shape, size, and color. It is usually several centimeters long and bright purple to pale green. It contains up to four seeds. Depending on the cultivar, the seeds are white, brown, red, or black, sometimes with a white hilum. Wild plants have mottled seeds. The seed is about a centimeter long.
Origin and occurrence
The exact origin of the lablab bean remains uncertain. Evidence of wild varieties in eastern and southern Africa suggests these regions as the likely source, although some theories suggest India as the origin.
Over the centuries, the lablab has been distributed all over the world. Despite its preference for tropical and subtropical climates, it can be found in temperate climates such as Central and South America or Italy. Its adaptability to different climates increases its agricultural value.
Agronomy
The lablab bean remains most widespread in tropical and subtropical areas, particularly in eastern and southern Africa and India. There, the legume is grown primarily for food and fodder, but its cultivation has declined sharply in many regions of Africa, despite renewed interest in its soil-improving functions in multiple cropping systems. One of the main reasons for this decline is the replacement of faba bean by the common bean, but factors such as the many processing steps and antinutritional factors may also have played a role. In addition, lablab is often grown in home garden systems or mixed cropping systems, such as in Southeast Asia, making estimates of global production difficult.
Growing requirements
Lablab can grow on a wide range of soils, from sand to clay, within a pH range of 4.5 to 7.5. It is known to grow better in acidic conditions than most legumes, but does not grow well in poorly drained soils or in saline conditions. Average daily temperatures of 18-30°C and annual rainfall of 700-3000 mm allow lablab cultivation. However, it can tolerate temperatures as low as 4°C for short periods. The seedlings grow slowly, but once established they compete well with weeds and are able to tolerate drought or shade. Lablab requires well-drained soils as it is very intolerant of waterlogged or flooded conditions.
Cultivation, harvest and storage
As there is a high degree of variability in phenotypes, including relative maturity, yield, susceptibility to insect attack and drought resistance, cultivation can vary accordingly depending on the accession and environmental factors However, when lablab is first planted in a field, it is beneficial to inoculate the seeds with specialised rhizobium bacteria. The seeds are either broadcast and then covered, or sown to a depth of 5 cm using the drillseed technique. The seedbed must be kept free of weeds in the early stages of growth, as the young plant succumbs easily to weed pressure.
It is generally accepted that harvesting fodder crops at the transition between the vegetative and reproductive stages gives the best compromise between yield and quality, as after this stage they become higher in fibre and lignin content and lower in protein content, leading to reduced digestibility and acceptability to livestock.
When lablab is stored as hay, the main challenge is the loss of dry matter from the leaves. Farmers suggest grinding it and storing it in bags. This can also prevent damage from direct sunlight or rain.
Soil improvements
Due to its symbiosis with nitrogen-fixing rhizobium bacteria, lablab bean has low soil fertility requirements and can be an important part of an agricultural system that improves soil nitrogen availability. The rate of nitrogen fixation has been shown to be highest when combined with phosphorus fertilisation, as phosphorus is required for rhizobium attachment to roots. Its cultivation also increases potassium and phosphorus availability. These soil-improving properties make lablab attractive for intercropping, mixed cropping systems and as a green manure. For example, intercropping lablab maize can improve multiple soil functions, including microbial diversity, with minimal to no loss of maize yield.
Breeding
Researchers suggest that the wide diversity of lablab germplasm has great potential for advancing the species as a promising alternative crop through selection and breeding.
Currently, most of the traditional cultivars grown have indeterminate growth habit, which has allowed farmers to harvest plants continuously.As such varieties aren't useful in the context of modern industrial farming, there has been a push to breed varieties that develop their seeds simultaneously, allowing all the beans to be collected in a single harvest. A side-effect of these breeding efforts has been a shortened vegetative phase, which reduces yield and hence economic value Breeding objectives also differ between countries and between uses of the plants. In Australia, the focus is on creating a variety for fodder use by crossing pure line varieties with wild varieties from Africa, while researchers in Bangladesh and Africa are trying to increase the yield of existing varieties to increase their value in the food production system.
Uses
The hyacinth bean is an old domesticated pulse and multi-purpose crop. L. purpureus has been cultivated in India as early as 2500 BC.
Due to seed availability of one forage cultivar (cv. Rongai), it is often grown as forage for livestock and as an ornamental plant. In addition, it is cited both as a medicinal plant and a poisonous plant.
The fruit and beans are edible if boiled well with several changes of the water. Otherwise, they are toxic due to the presence of cyanogenic glycosides, glycosides that are converted to hydrogen cyanide when consumed. Signs of poisoning include weakness, vomiting, shortness of breath, twitching, stupor, and convulsions. It has been shown that there is a wide range of cyanogenic potential among the varieties.
The leaves are eaten raw or cooked like spinach. The flowers can be eaten raw or steamed. The root can be boiled or baked for food. The seeds are used to make tofu and tempeh.
Food in South Asia
In Bangladesh and West Bengal, the green pods along with the beans, known as sheem (শিম), are cooked as vegetables or cooked with fish as a curry.
In Gujarat, lablab is called . In Maharashtra, they are known as val papdi. In both states, they are commonly used to make sautéed, spiced vegetable dishes called sabjis.
In Kerala, it is known as , or (Malayalam: ). The beans as well as the bean pods are used in cooking curries. The bean pods are also used (along with spices) for preparing a stir-fried dish known as .
In Tamil Nadu, it is called avarai or avaraikkaay (Tamil: அவரைக்காய் / அவரை). The entire bean is used in cooking dry curries and in sauces/gravies such as sambar. The seed alone is used in many recipes and is referred to as mochai (Tamil: மொச்சை / மொச்சைக்கொட்டை).
In Maharashtra, dry preparations with green masala are often made out of these green beans (ghevda varieties; Shravan ghevda (French beans), bajirao ghevda, ghevda, walwar, pavta sheng) mostly at the end of monsoon season during fasting festivals of Shravan month.
In Karnataka, the hyacinth bean is made into curry (avarekalu saaru) (), salad (avarekaalu usli), added to upma (avrekaalu uppittu), and as a flavoring to Akki rotti. Sometimes the outer peel of the seed is removed and the inner soft part is used for a variety of dishes. This form is called hitakubele avarekalu, which means "pressed (hitaku) hyacinth bean," and a curry known as hitikida avarekaalu saaru is made out of the deskinned beans.
In Telangana and Andhra Pradesh, the bean pods are cut into small pieces and cooked as a spicy curry in the Pongal festival season. Sometimes the outer peel of the seed when tender and soaked overnight is removed and the inner soft part is used for a variety of dishes. This form is called pitakapappu hanupa/anapa, which means "pressed (pitaku) hyacinth bean, and a curry known as pitikina anapaginjala chaaru/pitaka pappu is made from the deskinned beans and eaten along with bajra bread.
Food in Southeast and East Asia
In Myanmar, lablab beans are used to make a braised Burmese curry hnat (). They are also crisp-fried and served in Burmese pickled tea leaf salad.
In Huế, Vietnam, hyacinth beans are the main ingredient of the dish chè đậu ván (Hyacinth Bean Sweet Soup).
In China, the seeds are known as Bai Bian Dou. They are usually dried and baked before being used in traditional Chinese herbal remedies to strengthen the spleen, reduce heat and dampness, and promote appetite.
Food tradition in East Africa
In Kenya, the bean, known as njahe or njahi, is popular among several communities, especially the Kikuyu. Seasons were actually based on it, i.e., the Season of Njahe (Kīmera kīa njahī). It is thought to encourage lactation and has historically been the main dish for breastfeeding mothers. Beans are boiled and mashed with ripe and/or semi-ripe bananas, giving the dish a sweet taste. Today the production is in decline in eastern Africa. This is partly attributed to the fact that under colonial rule in Kenya, farmers were forced to give up their local bean in order to produce common beans (Phaseolus vulgaris) for export.
Medicinal use
Taiwanese research found that a carbohydrate-binding protein (i.e. a legume lectin) from lablab beans effectively blocks the infections of influenza viruses and SARS-CoV-2.
Gallery
| Biology and health sciences | Pulses | Plants |
1452308 | https://en.wikipedia.org/wiki/Borosilicate%20glass | Borosilicate glass | Borosilicate glass is a type of glass with silica and boron trioxide as the main glass-forming constituents. Borosilicate glasses are known for having very low coefficients of thermal expansion (≈3 × 10−6 K−1 at 20 °C), making them more resistant to thermal shock than any other common glass. Such glass is subjected to less thermal stress and can withstand temperature differentials without fracturing of about . It is commonly used for the construction of reagent bottles and flasks, as well as lighting, electronics, and cookware. For many other applications, soda-lime glass is more common.
Borosilicate glass is sold under various trade names, including Borosil, Duran, Pyrex, Glassco, Supertek, Suprax, Simax, Bellco, Marinex (Brazil), BSA 60, BSC 51 (by NIPRO), Heatex, Endural, Schott, Refmex, Kimax, Gemstone Well, United Scientific, and MG (India).
Single-ended self-starting lamps are insulated with a mica disc and contained in a borosilicate glass gas discharge tube (arc tube) and a metal cap. They include the sodium-vapor lamp that is commonly used in street lighting.
Borosilicate glass usually melts at about .
History
Borosilicate glass was first developed by German glassmaker Otto Schott in the late 19th century in Jena. This early borosilicate glass thus came to be known as Jena glass. After Corning Glass Works introduced Pyrex in 1915, the name became synonymous with borosilicate glass in the English-speaking world (since the 1940s, a sizable portion of glass produced under the Pyrex brand has also been made of soda–lime glass). Borosilicate glass is the name of a glass family with various members tailored to completely different purposes. Most common today is borosilicate 3.3 or 5.0x glass such as Duran, Corning33, Corning51-V (clear), Corning51-L (amber), International Cookware's NIPRO BSA 60, and BSC 51.
Manufacturing process
Borosilicate glass is created by combining and melting boric oxide, silica sand, soda ash, and alumina. Since borosilicate glass melts at a higher temperature than ordinary silicate glass, some new techniques were required for industrial production.
In addition to quartz, sodium carbonate, and aluminium oxide traditionally used in glassmaking, boron is used in the manufacture of borosilicate glass. The composition of low-expansion borosilicate glass, such as those laboratory glasses mentioned above, is approximately 80% silica, 13% boric oxide, 4% sodium oxide or potassium oxide and 2–3% aluminium oxide. Though more difficult to make than traditional glass due to its high melting temperature, it is economical to produce. Its superior durability, chemical and heat resistance finds use in chemical laboratory equipment, cookware, lighting, and in certain kinds of windows.
The manufacturing process depends on the product geometry and can be differentiated between different methods like floating, tube drawing, or molding.
Physical characteristics
The common type of borosilicate glass used for laboratory glassware has a very low thermal expansion coefficient (3.3 × 10−6 K−1), about one-third that of ordinary soda–lime glass. This reduces material stresses caused by temperature gradients, which makes borosilicate a more suitable type of glass for certain applications (see below). Fused quartzware is even better in this respect (having one-fifteenth the thermal expansion of soda–lime glass); however, the difficulty of working with fused quartz makes quartzware much more expensive, and borosilicate glass is a low-cost compromise. While more resistant to thermal shock than other types of glass, borosilicate glass can still crack or shatter when subjected to rapid or uneven temperature variations.
Among the characteristic properties of this glass family are:
Different borosilicate glasses cover a wide range of different thermal expansions, enabling direct seals with various metals and alloys like molybdenum glass with a CTE (coefficient of thermal expansion) of 4.6, tungsten with a CTE around 4.0 and Kovar with a CTE around 5.0 because of the matched CTE with the sealing partner
Allowing high maximum temperatures of typically about
Showing an extremely high chemical resistance in corrosive environments. Norm tests for example for acid resistance create extreme conditions and reveal very low impacts on glass
The softening point (temperature at which viscosity is approximately 107.6 poise) of type 7740 Pyrex is .
Borosilicate glass is less dense (about 2.23 g/cm3) than typical soda–lime glass due to the low atomic mass of boron. Its mean specific heat capacity at constant pressure (20–100 °C) is 0.83 J/(g⋅K), roughly one fifth of water's.
The temperature differential that borosilicate glass can withstand before fracturing is about , whereas soda–lime glass can withstand only about a change in temperature. This is why typical kitchenware made from traditional soda–lime glass will shatter if a vessel containing boiling water is placed on ice, but Pyrex or other borosilicate laboratory glass will not.
Optically, borosilicate glasses are crown glasses with low dispersion (Abbe numbers around 65) and relatively low refractive indices (1.51–1.54 across the visible range).
Families
For the purposes of classification, borosilicate glass can be roughly arranged in the following groups, according to their oxide composition (in mass fractions). Characteristic of borosilicate glasses is the presence of substantial amounts of silica (SiO2) and boric oxide (B2O3, >8%) as glass network formers. The amount of boric oxide affects the glass properties in a particular way. Apart from the highly resistant varieties (B2O3 up to a maximum of 13%), there are others that – due to the different way in which the boric oxide is incorporated into the structural network – have only low chemical resistance (B2O3 content over 15%). Hence we differentiate between the following subtypes.
Non-alkaline-earth
The B2O3 content for borosilicate glass is typically 12–13% and the SiO2 content over 80%. High chemical durability and low thermal expansion (3.3 × 10−6 K−1) – the lowest of all commercial glasses for large-scale technical applications – make this a versatile glass material. High-grade borosilicate flat glasses are used in a wide variety of industries, mainly for technical applications that require either good thermal resistance, excellent chemical durability, or high light transmission in combination with a pristine surface quality. Other typical applications for different forms of borosilicate glass include glass tubing, glass piping, glass containers, etc. especially for the chemical industry.
Alkaline-earth
In addition to about 75% SiO2 and 8–12% B2O3, these glasses contain up to 5% oxides of alkaline earth metal and alumina (Al2O3). This is a subtype of slightly softer glasses, which have thermal expansions in the range (4.0–5.0) × 10−6 K−1.
This is not to be confused with simple borosilicate glass-alumina composites.
High-borate
Glasses containing 15–25% B2O3, 65–70% SiO2, and smaller amounts of alkalis and Al2O3 as additional components have low softening points and low thermal expansion. Sealability to metals in the expansion range of tungsten and molybdenum and high electrical insulation are their most important features. The increased B2O3 content reduces the chemical resistance; in this respect, high-borate borosilicate glasses differ widely from non-alkaline-earth and alkaline-earth borosilicate glasses. Among these are also borosilicate glasses that transmit UV light down to 180 nm, which combine the best of the borosilicate glass and the quartz world.
Uses
Borosilicate glass has a wide variety of uses ranging from cookware to lab equipment, as well as a component of high-quality products such as implantable medical devices and devices used in space exploration.
Health and science
Virtually all modern laboratory glassware is made of borosilicate glass. It is widely used in this application due to its chemical and thermal resistance and good optical clarity, but the glass can react with sodium hydride upon heating to produce sodium borohydride, a common laboratory reducing agent. Fused quartz is also found in some laboratory equipment when its higher melting point and transmission of UV are required (e.g. for tube furnace liners and UV cuvettes), but the cost and manufacturing difficulties associated with fused quartz make it an impractical investment for the majority of laboratory equipment.
Additionally, borosilicate tubing is used as the feedstock for the production of parenteral drug packaging, such as vials and pre-filled syringes, as well as ampoules and dental cartridges. The chemical resistance of borosilicate glass minimizes the migration of sodium ions from the glass matrix, thus making it well suited for injectable-drug applications. This type of glass is typically referred to as USP / EP JP Type I.
Borosilicate is widely used in implantable medical devices such as prosthetic eyes, artificial hip joints, bone cements, dental composite materials (white fillings) and even in breast implants.
Many implantable devices benefit from the unique advantages of borosilicate glass encapsulation. Applications include veterinary tracking devices, neurostimulators for the treatment of epilepsy, implantable drug pumps, cochlear implants, and physiological sensors.
Electronics
During the mid-20th century, borosilicate glass tubing was used to pipe coolants (often distilled water) through high-power vacuum-tube–based electronic equipment, such as commercial broadcast transmitters. It was also used for the envelope material for glass transmitting tubes which operated at high temperatures.
Borosilicate glasses also have an application in the semiconductor industry in the development of microelectromechanical systems (MEMS), as part of stacks of etched silicon wafers bonded to the etched borosilicate glass.
Cookware
Cookware is another common usage for borosilicate glass, including bakeware. It is used for some measuring cups, featuring screen printed markings providing graduated measurements. Borosilicate glass is sometimes used for high-quality beverage glassware, particularly in pieces designed for hot drinks. Items made of borosilicate glass can be thin yet durable, or thicker for extra strength, and are microwave- and dishwasher-safe.
Lighting
Many high-quality flashlights use borosilicate glass for the lens. This increases light transmittance through the lens compared to plastics and lower-quality glass.
Several types of high-intensity discharge (HID) lamps, such as mercury-vapor and metal-halide lamps, use borosilicate glass as the outer envelope material.
New lampworking techniques led to artistic applications such as contemporary glass marbles. The modern studio glass movement has responded to color. Borosilicate is commonly used in the glassblowing form of lampworking and the artists create a range of products such as jewelry, kitchenware, sculpture, as well as for artistic glass smoking pipes.
Lighting manufacturers use borosilicate glass in some of their lenses.
Organic light-emitting diodes (OLED) (for display and lighting purposes) also use borosilicate glass (BK7). The thicknesses of the BK7 glass substrates are usually less than 1 millimeter for OLED fabrication. Due to its optical and mechanical characteristics in relation with cost, BK7 is a common substrate in OLEDs. However, depending on the application, soda–lime glass substrates of similar thicknesses are also used in OLED fabrication.
Optics
Many astronomical reflecting telescopes use glass mirror components made of borosilicate glass because of its low coefficient of thermal expansion. This makes very precise optical surfaces possible that change very little with temperature, and matched glass mirror components that "track" across temperature changes and retain the optical system's characteristics.
The Hale Telescope's 200 inch mirror is made of borosilicate glass.
The optical glass most often used for making instrument lenses is Schott BK-7 (or the equivalent from other makers, such as the Chinese crown glass K9), a very finely made borosilicate crown glass.
It is also designated as 517642 glass after its 1.517 refractive index and 64.2 Abbe number. Other less costly borosilicate glasses, such as Schott B270 or the equivalent, are used to make "crown-glass" eyeglass lenses. Ordinary lower-cost borosilicate glass, like that used to make kitchenware and even reflecting telescope mirrors, cannot be used for high-quality lenses because of the striations and inclusions common to lower grades of this type of glass. The maximal working temperature is . While it transitions to a liquid starting at (just before it turns red-hot), it is not workable until it reaches over . That means that in order to industrially produce this glass, oxygen/fuel torches must be used. Glassblowers borrowed technology and techniques from welders.
Rapid prototyping
Borosilicate glass has become the material of choice for fused deposition modeling (FDM), or fused filament fabrication (FFF), build plates. Its low coefficient of expansion makes borosilicate glass, when used in combination with resistance-heating plates and pads, an ideal material for the heated build platform onto which plastic materials are extruded one layer at a time. The initial layer of build must be placed onto a substantially flat, heated surface to minimize shrinkage of some build materials (ABS, polycarbonate, polyamide, etc.) due to cooling after deposition. Depending on the material used, the build plate will cycle from room temperature to between 50 °C and 130 °C for each prototype that is built. The temperature, along with various coatings (Kapton tape, painter's tape, hair spray, glue stick, ABS+acetone slurry, etc.), ensure that the first layer may be adhered to and remain adhered to the plate, without warping, as the first and subsequent layers cool following extrusion. Subsequently, following the build, the heating elements and plate are allowed to cool. The resulting residual stress formed when the plastic contracts as it cools, while the glass remains relatively dimensionally unchanged due to the low coefficient of thermal expansion, provides a convenient aid in removing the otherwise mechanically bonded plastic from the build plate. In some cases the parts self-separate as the developed stresses overcome the adhesive bond of the build material to the coating material and underlying plate.
Other
Aquarium heaters are sometimes made of borosilicate glass. Due to its high heat resistance, it can tolerate the significant temperature difference between the water and the nichrome heating element.
Specialty glass smoking pipes for cannabis and tobacco can be made from borosilicate glass. The high heat resistance makes the pipes more durable. Some harm reduction organizations also give out borosilicate pipes intended for smoking crack cocaine, as the heat resistance prevents the glass from cracking, causing cuts and burns that can spread hepatitis C.
Most premanufactured glass guitar slides are made of borosilicate glass.
Borosilicate is also a material of choice for evacuated-tube solar thermal technology because of its high strength and heat resistance.
The thermal insulation tiles on the Space Shuttle were coated with a borosilicate glass.
Borosilicate glasses are used for immobilisation and disposal of radioactive wastes. In most countries high-level radioactive waste has been incorporated into alkali borosilicate or phosphate vitreous waste forms for many years; vitrification is an established technology. Vitrification is a particularly attractive immobilization route because of the high chemical durability of the vitrified glass product. The chemical resistance of glass can allow it to remain in a corrosive environment for many thousands or even millions of years.
Borosilicate glass tubing is used in specialty TIG welding torch nozzles in place of standard alumina nozzles. This allows a clear view of the arc in situations where visibility is limited.
Trade names
Borosilicate glass is offered in slightly different compositions under different trade names:
Borofloat of Schott AG, a borosilicate glass, which is produced to flat glass in a float process.
Borosil, manufactured by the company of the same name, used in laboratory glassware and microwaveable kitchenware in India
BK7 of Schott, a borosilicate glass with a high level of purity. Main use in lens and mirrors for laser, cameras and telescopes.
Duran of DURAN Group, similar to Pyrex, Simax or Jenaer Glas.
Pyrex borosilicate glass of Corning
Fiolax of Schott, mainly used for containers for pharmaceutical applications.
Ilmabor of (2014 insolvency), mainly used for containers and equipment in laboratories and medicine.
Jenaer Glas of Zwiesel Kristallglas, formerly Schott AG. Mainly used for kitchenware.
Kimax is the trademark for borosilicate glassware from Kimble
United Scientific, manufacturers and distributors of laboratory glassware
Rasotherm of VEB Jenaer Glaswerk Schott & Genossen, for technical glass
Simax of Kavalierglass a.s., Czechia, produced for both laboratory and consumer markets.
Supertek, manufacturer of scientific lab equipment and glassware.
Willow Glass is an alkali free, thin and flexible borosilicate glass of Corning
Boroux is a brand of borosilicate glass drinking bottles.
Endural is a brand name of Holophane
Borosilicate nanoparticles
It was initially thought that borosilicate glass could not be formed into nanoparticles, since an unstable boron oxide precursor would prevent successful forming of these shapes. However, in 2008 a team of researchers from the Swiss Federal Institute of Technology at Lausanne were successful in forming borosilicate nanoparticles of 100 to 500 nanometers in diameter. The researchers formed a gel of tetraethylorthosilicate and trimethoxyboroxine. When this gel is exposed to water under proper conditions, a dynamic reaction ensues, which results in the nanoparticles.
In lampworking
Borosilicate (or "boro", as it is often called) is used extensively in the glassblowing process lampworking; the glassworker uses a burner torch to melt and form glass, using a variety of metal and graphite tools to shape it. Borosilicate is referred to as "hard glass" and has a higher melting point (approximately 3,000 °F / 1648 °C) than "soft glass", which is preferred for glassblowing by beadmakers. Raw glass used in lampworking comes in glass rods for solid work and glass tubes for hollow work tubes and vessels/containers. Lampworking is used to make complex and custom scientific apparatus; most major universities have a lampworking shop to manufacture and repair their glassware. For this kind of "scientific glassblowing", the specifications must be exact and the glassblower must be highly skilled and able to work with precision. Lampworking is also done as art, and common items made include goblets, paper weights, pipes, pendants, compositions and figurines.
In 1968, English metallurgist John Burton brought his hobby of hand-mixing metallic oxides into borosilicate glass to Los Angeles. Burton began a glass workshop at Pepperdine College, with instructor Margaret Youd. A few of the students in the classes, including Suellen Fowler, discovered that a specific combination of oxides made a glass that would shift from amber to purples and blues, depending on the heat and flame atmosphere. Fowler shared this combination with Paul Trautman, who formulated the first small-batch colored borosilicate recipes. He then founded Northstar Glassworks in the mid-1980s, the first factory devoted solely to producing colored borosilicate glass rods and tubes for use by artists in the flame. Trautman also developed the techniques and technology to make the small-batch colored boro that is used by a number of similar companies.
Beadmaking
In recent years, with the resurgence of lampworking as a technique to make handmade glass beads, borosilicate has become a popular material in many glass artists' studios. Borosilicate for beadmaking comes in thin, pencil-like rods. Glass Alchemy, Trautman Art Glass, and Northstar are popular manufacturers, although there are other brands available. The metals used to color borosilicate glass, particularly silver, often create strikingly beautiful and unpredictable results when melted in an oxygen-gas torch flame. Because it is more shock-resistant and stronger than soft glass, borosilicate is particularly suited for pipe making, as well as sculpting figures and creating large beads. The tools used for making glass beads from borosilicate glass are the same as those used for making glass beads from soft glass.
| Technology | Materials | null |
1452399 | https://en.wikipedia.org/wiki/Vigna%20aconitifolia | Vigna aconitifolia | Vigna aconitifolia is a drought-resistant legume, commonly grown in arid and semi-arid regions of India. It is commonly called mat bean, moth bean, matki or dew bean. The pods, sprouts and protein-rich seeds of this crop are commonly consumed in India. Moth bean can be grown on many soil types, and can also act as a pasture legume.
Moth bean is a creeping annual herbaceous plant which grows to approximately 40 cm high. Yellow flowers on its hairy and densely packed branches develop into yellow-brown pods, 2 to 3 inches in length The seeds of these pods contain approximately 22–24% protein.
Due to its drought-resistant qualities, its ability to combat soil erosion and its high protein content, moth bean has been identified as possibly a more significant food source in the future. It has been suggested that its suitability as a grain legume in semi-arid Africa should be further investigated.
Description
Taxonomically moth bean belongs to the family Fabaceae (subfamily Papilionoideae),and genus Vigna. It is an herbaceous creeping annual that creates a low-lying soil cover when fully grown. Its stem can grow up to 40 cm in height, with its hairy and dense-packed branches reaching a span of up to 150 cm. Yellow flowers develop into a brown pod 2.5 to 5 cm in length, which holds 4 to 9 seeds inside. The rectangular seeds exist in a variety of colours including yellow-brown, whitish-green, and mottled with black. Other widely cultivated species from the genus Vigna include the adzuki bean (V. angularis), the black gram (V. mungo), the cowpea (V. unguiculata, including the variety known as the black-eyed pea), and the mung bean (V. radiata).
History and geography
Moth bean is native to India, grown for food production and as a forage and cover crop. It is predominately grown in India, although it has been cultivated in the United States, Australia, Thailand and other parts of Asia. 1.5 million hectares of land is used in India for moth bean production, producing approximately 0.4 million t/ha of seeds. While its presence in Sudan, Somalia and other tropical countries of Africa has been noted, it has not been a crop of great importance to this region. The potential of increased production in this region in the future has been suggested.
Growing conditions
Moth bean, a short-day crop, is one of the most drought-resistant pulses in India. Grown at altitudes up to 1300 m above sea level, it has a wide pH range (3.5–10) and can tolerate slight salinity. While dry sandy soil is most suitable for production, moth bean can tolerate a variety of soil types. The low-lying soil cover the crop creates helps prevent soil erosion by preventing moisture loss.
Optimum production of moth bean occurs between 24–32 °C, but has been shown to tolerate up to 45 °C during the day. Growth is optimal at a constant temperature. The moth bean is one of the most drought-resistant pulses in India, requiring little irrigation for production. While optimal annual rainfall for production is 500–750 mm, it is able to grow with 200–300 mm annually, and some yield has been noted at rainfall levels as low as 50–60 mm per year. Propagation of moth bean is done by seed, preferably on a prepared seedbed, at an optimal temperature of 25–27 °C. Fertilizer applications to moth bean are uncommon in India.
Other farming issues
Moth bean is grown for both human consumption and as a forage crop. Currently in India, moth bean is grown on its own or intercropped with other cereals, such as pearl millet. It is also grown in rotation with cotton as a forage crop.
When grown as a forage crop, it is planted 7–34 kg/ha, and 10–20 kg/ha when grown as the only crop. Row planting should be done 30–90 cm apart, with seeds sown 2.5–4 cm deep. It takes 75–90 days for moth bean to mature, and is frequently planted at the end of the rainy season.
A drawback to this crop is its difficulty to harvest. Mowers cannot be used due to the shape and density of the moth bean's branches, so the crop is typically cut with a sickle. It is threshed and winnowed after being dried for approximately one week. Hay yields from this crop are 7.5-10 t/ha, while forage matter yields range from 37-50 t/ha. Seed yields are currently low, ranging from 70–270 kg/ha. However, research shows that this crop has the potential to increase in yield. Experimental seed yields of up to 2600 kg/ha have been recorded in the US and Australia.
Major pests and diseases
Moth bean is affected by mungbean yellow mosaic virus, for which silverleaf whitefly is the vector. Root rot and seedling blight from Macrophomina phaseolina also cause damage, as well as some Striga species and the nematode Meloidogyne incognita. There are some resistant cultivars to these pests and diseases.
Genetic stock
Little breeding work has been completed on the moth bean, but researchers have found that there is substantial genetic variation between moth bean germplasms. The National Bureau of Plant Genetic Resources in New Delhi, India, houses more than 1000 accessions. Some improved cultivars such as ' CZM-2, CZM-3, ‘RMO-40’ and ‘RMO-225’are available in India.
Human consumption
Whole or split moth bean seeds can be cooked or fried. In India, particularly in the state of Maharashtra, moth beans are sprouted before cooking and used for making a spicy stew called matki usal. Matki usal forms the base of the Indian street food dish called misal pav. "Fried dal of the bean is used for making a savory dry snack, in India called dalmoth. The moth bean pods can be boiled and eaten. The flour of the bean is used for making another savoury snack called bhujia."
Animal consumption
Moth bean is also consumed as a forage crop by animals.
Nutritional information
100g of raw, uncooked moth bean seeds contain 343 calories, 23 g of protein, 62 g of carbohydrate and 1.6 g of fat. As is the case with other legumes, this pulse does contain antinutritional factors that limit available protein. However, research has shown that the moth bean contains considerably less of these factors compared with other legume grains, making it a more beneficial choice for consumption. Soaking and cooking moth beans before consumption helps to break down antinutritional factors and makes the protein more digestible.
Constraints to wider adaptation
While its drought tolerance and high protein content could make the moth bean a potential crop choice for semi-arid Africa, a lack of management knowledge and the difficulty of harvest due to its density and creeping nature could make its spread to other parts of the world difficult.
| Biology and health sciences | Pulses | Plants |
1454654 | https://en.wikipedia.org/wiki/Spinning%20frame | Spinning frame | The spinning frame is an Industrial Revolution invention for spinning thread or yarn from fibres such as wool or cotton in a mechanized way. It was developed in 18th-century Britain by Richard Arkwright and John Kay.
Historical context
In 1760 England, yarn production from wool, flax and cotton was still a cottage industry in which fibres were carded and spun by hand using a spinning wheel. As the textile industry expanded its markets and adopted faster machines, yarn supplies became scarce especially due to innovations such as the doubling of the loom speed after the invention of the flying shuttle. High demand for yarn spurred invention of the spinning jenny in 1764, followed closely by the invention of the spinning frame, later developed into the water frame (patented in 1769). Mechanisms had increased production of yarn so dramatically that by 1830 the cottage yarn industry in England could no longer compete and all spinning was carried out in factories.
Development
Richard Arkwright employed John Kay to produce a new spinning machine that Kay had worked on with (or possibly stolen from) another inventor named Thomas Highs. With the help of other local craftsmen, including Peter Atherton, the team developed the spinning frame, which produced a stronger thread than the spinning jenny invented by James Hargreaves. The frame utilised the draw rollers invented by Lewis Paul to stretch, or attenuate, the yarn.
The roller spinning process starts with a thick 'string' of loose fibres called a roving, which is passed between three pairs of rollers, each pair rotating slightly faster than the previous one. In this way it is reduced in thickness and increased in length before a strengthening twist is added by a bobbin-and-flyer mechanism. The spacing of the rollers has to be slightly greater than the fiber length to prevent breakage. The nip of the roller pairs prevents the twist from backing up to the roving.
Too large to be operated by hand, the spinning frame needed a new source of power. Arkwright experimented with horses, but decided to employ the power of the water wheel, which gave the invention the name 'water frame'.
For some time, the stronger yarn produced by the spinning frame was used in looms for the lengthwise warp threads that bound cloth together, while hand powered jennies provided the weaker yarn used for the horizontal filler (weft) threads. The jennies required skill but were inexpensive and could be used in a home. The spinning frames required capital but little skill.
In France, Philippe Henri de Girard patented spinning frames for both the dry and wet spinning of flax. His inventions were also patented in England in 1815, in the name of Horace Hall.
Little is known about Horace Hall, it is a possible pseudonym. Undoubtedly if he had taken out his patent in England in 1815, the year of the Battle of Waterloo, in a French name there would not have been a lot of enthusiasm for it; possibly even suspicion.
In the British Isles, James Kay was initially credited with the invention of this device. On 2 December 1826 shortly after Kay's patent was awarded, Philippe Henri de Girard wrote to the Editor of The Manchester Guardian complaining about and pointing out he had been the inventor. The following is an extract from his letter:
Kay's patent was invalidated, in 1839, on the grounds it was too similar to Horace Hall's; A decision upheld on appeal, in 1841.
| Technology | Spinning | null |
1454895 | https://en.wikipedia.org/wiki/Flying%20junction | Flying junction | A flying junction or flyover is a railway junction at which one or more diverging or converging tracks in a multiple-track route cross other tracks on the route by bridge to avoid conflict with other train movements. A more technical term is "grade-separated junction". A burrowing junction or dive-under occurs where the diverging line passes below the main line.
The alternative to grade separation is a level junction or flat junction, where tracks cross at grade, and conflicting routes must be protected by interlocked signals.
Complexity
Simple flying junctions may have a single track pass over or under other tracks to avoid conflict; complex flying junctions may have elaborate infrastructure to allow multiple routings without trains coming into conflict, in the manner of a highway stack interchange.
Flying junction without crossings
Where two lines each of two tracks merge with a flying junction, they can become a four-track railway together, the tracks paired by direction. This happens regularly in the Netherlands (see Examples below).
High-speed rail
Nearly all junctions with high-speed railways are grade-separated. On the French Lignes à Grande Vitesse (TGV) high-speed network, the principal junction on the LGV Sud-Est, at Pasilly where the line to Dijon diverges, and on the LGV Atlantique at Courtalain where the line to Le Mans diverges, are fully grade-separated with special high-speed switches (points in British terminology) that permit the normal line speed of on the main line, and a diverging speed of .
The LGV network has four grade-separated high-speed triangles: Fretin (near Lille), Coubert (southeast Paris), Claye-Souilly (northeast Paris) and Angles (Avignon). A fifth, Vémars (northeast Paris), is grade-separated except for a single-track link on the least-used side, linking Paris Gare du Nord and Paris CDG airport.
Examples
Australia
Bowen Hills railway station in Brisbane
Burnley railway station in Melbourne
Camberwell railway station in Melbourne
Sydney Central Station
Glenfield railway station, Sydney
Strathfield railway station
Sandgate Flyover, Newcastle – main line flies over coal branch line
Goodwood railway station in Adelaide
Bayswater railway station in Perth
Canada
Columbia station in New Westminster, BC – Expo Line branches for King George (top) and Production Way–University (bottom)
Bridgeport station in Richmond, BC – Canada Line branches for YVR–Airport and Richmond–Brighouse
Denmark
Hvidovre, Copenhagen ()
Junction of M1 and M2 lines on the Copenhagen Metro
Lunderskov ()
Roskilde, south of ()
Sydhavnen, Copenhagen ()
Vigerslev, Copenhagen ()
Finland
Railway junction of two main lines at Kytömaa, Kerava
France (LGV Triangles)
Triangle de Fretin, Lille, connecting Paris, Brussels and London
Triangle de Coubert, Paris
Triangle des Angles, Avignon, with two parallel viaducts
Triangle de Claye-Souilly, Paris, partial four-way junction
Triangle de Vémars, Paris
Germany
Bruchsal Rollenberg junction
Hong Kong
Where Airport Express and Tung Chung line diverge from each other at Tai Ho Wan
Tseung Kwan O line to the east of Tseung Kwan O station
Netherlands
There are between 25 and about 40 flying junctions on Dutch railways, depending on how more complex examples are counted.
Near Harmelen. Before conversion to a flying junction, this was the site of the Harmelen train disaster.
At Breukelen railway station
At Lage Zwaluwe railway station
Flying junctions where the merged lines become a four track railway:
Near Den Haag Laan van NOI railway station
North of Leiden where lines from Haarlem and Schiphol merge
At Boxtel railway station where lines from 's-Hertogenbosch and Tilburg merge
West of Gouda where lines from Rotterdam and The Hague merge
More complex flying junctions, with tracks from four directions joining:
Around Amsterdam Sloterdijk railway station
Around Duivendrecht railway station
Northwest exit of Utrecht Centraal railway station
West and northwest exit of Rotterdam Centraal railway station
At both sides of Weesp railway station (see diagram at right)
Norway
Lillestrøm ()
Lysaker ()
Sandvika, east of and west of () ()
Sweden
Flemingsberg ()
Järna, north of ()
Järna, south of ()
Lund ()
Hyllie ()
Myrbacken ()
Lernacken ()
Södertälje hamn ()
Södertälje syd ()
Tomteboda ()
Taiwan
Start of Shalun line, south of Zhongzhou railway station
United Kingdom
Pelaw Junction where both the Tyne and Wear Metro green line to South Hylton joins the Durham Coast Line and yellow line continues to South Shields – both diverging on the bridge itself
Springhead Junction on the North Kent Line
Southfleet Junction on the HS1
Norton Bridge Junction near Stone, Staffordshire
Hamilton Square underground station, Birkenhead, on Merseyrail
Aynho Junction in Aynho, Northamptonshire
Worting Junction near Basingstoke, Hampshire (the flyover is called Battledown Flyover)
Cogload Junction near Taunton
Weaver Junction near Dutton, Cheshire
Shortlands Junction in south London
Northwest of Harrow-on-the-Hill, in the north London suburbs
Hitchin flyover, Hertfordshire.
Werrington Junction dive-under, northern suburbs of Peterborough, Cambridgeshire
Reading West Junction
Bleach Green Viaducts & Junction, Whiteabbey, Northern Ireland
Manchester Metrolink, Greater Manchester, immediately southwest of Cornbrook tram stop where the Eccles Line diverges from the Altrincham Line.
United States
Northeast U.S. (Amtrak)
Along the New York–Washington section of the Northeast Corridor, and on the Philadelphia–Harrisburg section of the Keystone Corridor. Both converge at Zoo Junction near 30th Street Station in Philadelphia. All were built by the former Pennsylvania Railroad and are now maintained by Amtrak.
Boston, Massachusetts area
An abandoned underground junction on the Tremont Street subway approaching the Pleasant Street incline
The B branch splits off from the C and D branches of the MBTA Green Line via an underground flying junction just west of Kenmore station.
The Union Square spur splits off from the Green Line Extension of the MBTA Green Line via an aerial flying junction on the Red Bridge viaduct in the Inner Belt area of Somerville, Massachusetts just north of Lechmere station in Cambridge. Lead tracks to the GLX maintenance facility also split off from the junction slightly further west.
The two southern branches of the MBTA Red Line in Boston split via a flying junction just north of JFK/UMass station. In addition, lead tracks to Cabot Yard maintenance facilities branch off from the junction.
Chicago, Illinois
On the Chicago "L", where Orange Line trains diverge from Green Line trains north of 18th Street, as well as underground where a non-revenue flying junction separates Red Line trains heading to 95th from those heading to the South Side main line, currently used to send some rush-period Red Line trains to Ashland/63rd.
The Milwaukee–Dearborn subway (now part of the Blue Line) was constructed to have a flying junction where turning between Lake Street and Milwaukee Avenue at Canal Street. The outbound tunnel and its stub, designed to continue west under Lake Street, was bored at less depth than the inbound tunnel and its Lake Street stub, in order to allow future Lake Street trains (now part of the Green and (Pink Lines) to run under or over the opposing Milwaukee Avenue trains while entering or exiting the shared portion of the Lake Street tunnels. Plans in 1939 called for tunnels to replace the elevated Lake Street tracks east of approximately Racine Avenue. By 1962, the planned Lake Street tunnels to/from Racine Avenue would have curved south to Randolph Street and bypassed the Milwaukee-Lake-Dearborn tunnel entirely.
Another flying junction is under construction immediately north of Belmont/Sheffield to increase capacity on the Red Line, Brown Line, and Purple Line Express.
Denver, Colorado
On the Regional Transportation District in Denver between the Southeast Corridor and the I-225 Corridor: the Southeast Corridor is on the west side of I-25 and the I-225 Corridor is in the median of I-225. The grade separations of the junction are woven into the grade separations of the interchange between the two highways.
New York, New York
On the New York City Subway there is an above-ground example at Hammel's Wye on the IND Rockaway Line, as well as numerous below-ground examples across the network
Connecting Metro-North Railroad's New Haven Line and Harlem Line, near Wakefield station in the Bronx
Philadelphia, Pennsylvania
Amtrak's Zoo Junction, where the Northeast Corridor meets the Keystone Corridor and sorts into 30th Street Station's lower and upper level platforms. Also known as Zoo Interlocking, the name comes from the Philadelphia Zoo, which is located in the crescent shaped pocket between the junction and the river.
On SEPTA's Cynwyd Line, diverging from the Keystone Corridor west of 52nd Street.
On SEPTA Airport Line, diverging from the Northeast Corridor south of Penn Medicine
On SEPTA's subway–surface trolley lines, where the Route 10 diverges from Routes 11/13/34/36 west of 33rd Street.
On SEPTA's Broad Street subway, where Broad-Ridge Spur trains diverge at Fairmount. There are also provisions for flying junctions north of Erie for the Roosevelt Boulevard Subway, and north of Olney for an extension on North Broad Street; both are maintained as layup tracks.
San Francisco Bay Area, California
The Oakland Wye, where all of Bay Area Rapid Transit's mainline operations converge near downtown Oakland
On the Market Street subway in San Francisco where the J Church and N Judah lines join the main line of the subway. The subway portal is east of the intersection of Church Street and Duboce Avenue in the Duboce Triangle neighborhood, immediately north of a Safeway supermarket and south of the San Francisco branch of the United States Mint.
Washington, District of Columbia
All main-line connections on the Washington Metro – adjacent to the Pepco power plant on Benning Road (near the Stadium-Armory station) is a large three-track structure with a turnback pocket where the Blue, Silver and Orange Lines meet. This would have been part of the Oklahoma Avenue station, had it been built. South of the King Street station in Alexandria is a series of tunnels where the Blue and Yellow Lines meet. There are also flying junctions near three underground rail stations: Rosslyn (Blue, Silver, and Orange Lines), L'Enfant Plaza (Green and Yellow lines), and the Pentagon (Blue and Yellow lines).
| Technology | Transport infrastructure | null |
1455603 | https://en.wikipedia.org/wiki/BEST%20theorem | BEST theorem | In graph theory, a part of discrete mathematics, the BEST theorem gives a product formula for the number of Eulerian circuits in directed (oriented) graphs. The name is an acronym of the names of people who discovered it: N. G. de Bruijn, Tatyana Ehrenfest, Cedric Smith and W. T. Tutte.
Precise statement
Let G = (V, E) be a directed graph. An Eulerian circuit is a directed closed trail that visits each edge exactly once. In 1736, Euler showed that G has an Eulerian circuit if and only if G is connected and the indegree is equal to outdegree at every vertex. In this case G is called Eulerian. We denote the indegree of a vertex v by deg(v).
The BEST theorem states that the number ec(G) of Eulerian circuits in a connected Eulerian graph G is given by the formula
Here tw(G) is the number of arborescences, which are trees directed towards the root at a fixed vertex w in G. The number tw(G) can be computed as a determinant, by the version of the matrix tree theorem for directed graphs. It is a property of Eulerian graphs that tv(G) = tw(G) for every two vertices v and w in a connected Eulerian graph G.
Applications
The BEST theorem shows that the number of Eulerian circuits in directed graphs can be computed in polynomial time, a problem which is #P-complete for undirected graphs. It is also used in the asymptotic enumeration of Eulerian circuits of complete and complete bipartite graphs.
History
The BEST theorem is due to van Aardenne-Ehrenfest and de Bruijn
(1951), §6, Theorem 6.
Their proof is bijective and generalizes the de Bruijn sequences. In a "note added in proof", they refer to an earlier result by Smith and Tutte (1941) which proves the formula for graphs with deg(v)=2 at every vertex.
| Mathematics | Graph theory | null |
33198793 | https://en.wikipedia.org/wiki/Ligia | Ligia | Ligia is a genus of isopods, commonly known as rock lice or sea slaters. Most Ligia species live in tidal zone cliffs and rocky beaches, but there are several fully terrestrial species which occur in high-humidity environments.
Ecology
Coastal Ligia species exhibit a mixture of terrestrial and marine characteristics, drying out easily, needing moist air and proximity to water to retain water. While they have gills and can exchange gas under water, they only do so when escaping terrestrial predators or being dislodged by wave action. They do not move swiftly in the water and are open to marine predation. They are well adapted to rocky surfaces and avoid sand, which opens them to terrestrial predation and desiccation.
Taxonomy
It has been suggested that Ligia is more closely related to marine isopods than it is to true woodlice.
Species
Species separation is at times difficult because of sexual dimorphism. For example, males usually have longer and wider antennae than females. The males also tend to be larger but narrower, with the difference sometimes attributed to the female's brood pouch. Complicating matters is the possible existence of cryptic species in the genus.
This is a list of all Ligia species contained in A Bibliography of Terrestrial Isopods:
Ligia australiensis – Australian slater, Australia, including Tasmania and Lord Howe Island
Ligia baudiniana – east and west coasts of the Americas
Ligia boninensis – Bonin Islands, Japan
Ligia cajennensis – French Guiana
Ligia cinerascens – Hokkaido and northern Honshu, Japan, and the Kuril Islands
Ligia cursor – Chile
Ligia curvata – Angola
Ligia dante – Hawai'i Island (Hawaii)
Ligia dentipes – Andaman and Nicobar Islands, Indonesia
Ligia dilatata – southern Africa (Namibia and South Africa)
Ligia eleluensis – Maui, Hawaii
Ligia exotica – wharf roach, introduced around the world in subtropical and warm temperate coastlines
Ligia ferrarai – Madagascar
Ligia filicornis – Venezuela
Ligia glabrata – Southern Africa (Namibia and South Africa)
Ligia gracilipes – west coast of Africa, Senegal to northern Angola
Ligia hachijoensis – Izu Islands, Japan
Ligia hawaiensis – Kaua'i, Hawaii
Ligia honu – Hawaii Island, Hawaii
Ligia italica – coasts of the Black Sea, Mediterranean Sea, Atlantic in northern Africa downsouth to Cape Verde and Macaronesia Islands
Ligia kamehameha – Hawaii Island, Hawaii
Ligia latissima – New Caledonia
Ligia litigiosa – Chile and Peru and the Juan Fernández Islands
Ligia malleata – Tanzania
Ligia mauinuiensis – Maui Nui Islands and Oahu, Hawaii
Ligia miyakensis – Izu Islands, Japan
Ligia natalensis – southeastern coast of South Africa from Knysna to Natal
Ligia novaezealandiae – New Zealand and Kermadec Island
Ligia occidentalis – California, Baja California, British Columbia
Ligia oceanica – Atlantic Coasts of Europe, coasts of western Baltic Sea and possibly introduced to the Atlantic Coast of North America
Ligia pallasii – Pacific Coast of North America, the Aleutian Islands to Santa Cruz, California
Ligia pallida – Christmas Island in Polynesia
Ligia pele – Maui, Hawai'i
Ligia perkinsi – Kaua'i and O'ahu, Hawaiian Islands
Ligia persica – Persian Gulf
Ligia philoscoides – southeastern Polynesia
Ligia pigmentata – Red Sea, Persian Gulf and coast of Somalia
Ligia platycephala – Venezuela, Guyana and Trinidad
Ligia rolliensis – Oahu, Hawai'i
Ligia rugosa – southeastern Polynesia
Ligia ryukyuensis – Japan
Ligia saipanensi – Saipan Island, Micronesia
Ligia simoni – Northern Venezuela and northern Colombia
Ligia taiwanensis – Taiwan
Ligia vitiensis – Sulawesi, Singapore, New Guinea, Melanesia, Polynesia, and possibly introduced to Somalia
Ligia yamanishii – Tokyo
Ligia yemenica – Gulf of Aden
| Biology and health sciences | Malacostraca | Animals |
21580073 | https://en.wikipedia.org/wiki/Pollicipes%20pollicipes | Pollicipes pollicipes | Pollicipes pollicipes, known as the goose neck barnacle, goose barnacle or leaf barnacle is a species of goose barnacle, also well known under the taxonomic synonym Pollicipes cornucopia. It is closely related to Pollicipes polymerus, a species with the same common names, but found on the Pacific coast of North America, and to Pollicipes elegans a species from the coast of Chile. It is found on rocky shores in the north-east Atlantic Ocean and is prized as a delicacy, especially in the Iberian Peninsula.
Distribution
Pollicipes pollicipes is chiefly distributed from 48°N to 28°N, along the coasts of France, Spain (including the Canary Islands), Portugal, Morocco, and south to Senegal. The periphery of the species' range also extends as far north as Ireland, with outlying populations on the south coast of England and possibly in southwestern Ireland, although there are no recent records there. The species is present, but rare, in the Mediterranean Sea. It is possible that the outlying populations are not self-sustaining, being instead maintained by immigration of larvae from self-sustaining core populations.
A population disjunctly located around the tropical Cape Verde Islands at about 16°N was described in 2010 as a new species, Pollicipes caboverdensis.
Ecology
Pollicipes pollicipes grows in groups on rocks, as well as on the hulls of shipwrecks and on driftwood. It is a filter feeder, living on particles that it can glean from the water passing over its extended cirri; these possess a complex assortment of setae, enabling P. pollicipes to have a varied diet, including diatoms, detritus, large crustaceans, copepods, shrimp and molluscs.
The larvae pass through seven free-swimming stages (six nauplii and one cypris) over the course of at least a month. After this time, they settle into the adult, sessile form.
Pollicipes pollicipes is harvested for consumption in many parts of its range, mostly for the Spanish market, where (marketed as percebe gallego) it may sell for as much as €90 per kilogram. As a result, the species is thought to be in decline. It is harvested manually, and archaeological evidence suggests that the species has been harvested in this way for over 10,000 years.
| Biology and health sciences | Crustaceans | Animals |
21581881 | https://en.wikipedia.org/wiki/Pharynx | Pharynx | The pharynx (: pharynges) is the part of the throat behind the mouth and nasal cavity, and above the esophagus and trachea (the tubes going down to the stomach and the lungs respectively). It is found in vertebrates and invertebrates, though its structure varies across species. The pharynx carries food to the esophagus and air to the larynx. The flap of cartilage called the epiglottis stops food from entering the larynx.
In humans, the pharynx is part of the digestive system and the conducting zone of the respiratory system. (The conducting zone—which also includes the nostrils of the nose, the larynx, trachea, bronchi, and bronchioles—filters, warms and moistens air and conducts it into the lungs). The human pharynx is conventionally divided into three sections: the nasopharynx, oropharynx, and laryngopharynx.
In humans, two sets of pharyngeal muscles form the pharynx and determine the shape of its lumen. They are arranged as an inner layer of longitudinal muscles and an outer circular layer.
Structure
Nasopharynx
The upper portion of the pharynx, the nasopharynx, extends from the base of the skull to the upper surface of the soft palate. It includes the space between the internal nares and the soft palate and lies above the oral cavity. The adenoids, also known as the pharyngeal tonsils, are lymphoid tissue structures located in the posterior wall of the nasopharynx. Waldeyer's tonsillar ring is an annular arrangement of lymphoid tissue in both the nasopharynx and oropharynx. The nasopharynx is lined by respiratory epithelium that is pseudostratified, columnar, and ciliated.
Polyps or mucus can obstruct the nasopharynx, as can congestion due to an upper respiratory infection. The auditory tube, which connects the middle ear to the pharynx, opens into the nasopharynx at the pharyngeal opening of the auditory tube. The opening and closing of the auditory tubes serves to equalize the barometric pressure in the middle ear with that of the ambient atmosphere.
The anterior aspect of the nasopharynx communicates through the choanae with the nasal cavities. On its lateral wall is the pharyngeal opening of the auditory tube, somewhat triangular in shape and bounded behind by a firm prominence, the torus tubarius or cushion, caused by the medial end of the cartilage of the tube that elevates the mucous membrane.
Two folds arise from the cartilaginous opening:
the salpingopharyngeal fold, a vertical fold of mucous membrane extending from the inferior part of the torus and containing the salpingopharyngeus muscle.
the salpingopalatine fold, a smaller fold, in front of the salpingopharyngeal fold, extending from the superior part of the torus to the palate and containing the salpingopalatine muscle. The tensor veli palatini and levator veli palatini are lateral to the fold and do not contribute.
Oropharynx
The oropharynx lies behind the oral cavity, extending from the uvula to the level of the hyoid bone. It opens anteriorly, through the isthmus faucium, into the mouth, while in its lateral wall, between the palatoglossal arch and the palatopharyngeal arch, is the palatine tonsil. The anterior wall consists of the base of the tongue and the epiglottic vallecula; the lateral wall is made up of the tonsil, tonsillar fossa, and tonsillar (faucial) pillars; the superior wall consists of the inferior surface of the soft palate and the uvula. Because both food and air pass through the pharynx, a flap of connective tissue called the epiglottis closes over the glottis when food is swallowed to prevent aspiration. The oropharynx is lined by non-keratinized squamous stratified epithelium.
The HACEK organisms (Haemophilus, Actinobacillus actinomycetemcomitans, Cardiobacterium hominis, Eikenella corrodens, Kingella) are part of the normal oropharyngeal flora, which grow slowly, prefer a carbon dioxide-enriched atmosphere, and share an enhanced capacity to produce endocardial infections, especially in young children. Fusobacterium is a pathogen.
Laryngopharynx
The laryngopharynx, (Latin: pars laryngea pharyngis), also known as hypopharynx, is the caudal part of the pharynx; it is the part of the throat that connects to the esophagus. It lies inferior to the epiglottis and extends to the location where this common pathway diverges into the respiratory (laryngeal) and digestive (esophageal) pathways. At that point, the laryngopharynx is continuous with the esophagus posteriorly. The esophagus conducts food and fluids to the stomach; air enters the larynx anteriorly. During swallowing, food has the "right of way", and air passage temporarily stops. Corresponding roughly to the area located between the 4th and 6th cervical vertebrae, the superior boundary of the laryngopharynx is at the level of the hyoid bone. The laryngopharynx includes three major sites: the pyriform sinus, postcricoid area, and the posterior pharyngeal wall. Like the oropharynx above it, the laryngopharynx serves as a passageway for food and air and is lined with a stratified squamous epithelium. It is innervated by the pharyngeal plexus and by the recurrent laryngeal nerve.
The vascular supply to the laryngopharynx includes the superior thyroid artery, the lingual artery and the ascending pharyngeal artery. The primary neural supply is from both the vagus and glossopharyngeal nerves. The vagus nerve provides an auricular branch also termed "Arnold's nerve" which also supplies the external auditory canal, thus laryngopharyngeal cancer can result in referred ear pain. This nerve is also responsible for the ear-cough reflex in which stimulation of the ear canal results in a person coughing.
Function
The pharynx moves food from the mouth to the esophagus. It also moves air from the nasal and oral cavities to the larynx. It is also used in human speech, as pharyngeal consonants are articulated here, and it acts as a resonating chamber during phonation.
Clinical significance
Inflammation
Inflammation of the pharynx, or pharyngitis, is the painful inflammation of the throat.
Pharyngeal cancer
Pharyngeal cancer is a cancer that originates in the neck and/or throat.
Waldeyer's tonsillar ring
Waldeyer's tonsillar ring is an anatomical term collectively describing the annular arrangement of lymphoid tissue in the pharynx. Waldeyer's ring circumscribes the naso- and oropharynx, with some of its tonsillar tissue located above and some below the soft palate (and to the back of the oral cavity). It is believed that Waldeyer's ring prevents the invasion of microorganisms from going into the air and food passages and this helps in the defense mechanism of the respiratory and alimentary systems.
Etymology
The word pharynx () is derived from the Greek φάρυγξ phárynx, meaning "throat". Its plural form is pharynges or pharynxes , and its adjective form is pharyngeal ( or ).
Other vertebrates
All vertebrates have a pharynx, used in both feeding and respiration. The pharynx arises during development in all vertebrates through a series of six or more outpocketings on the lateral sides of the head. These outpocketings are pharyngeal arches, and they give rise to a number of different structures in the skeletal, muscular, and circulatory systems. The structure of the pharynx varies across the vertebrates. It differs in dogs, horses, and ruminants. In dogs, a single duct connects the nasopharynx to the nasal cavity. The tonsils are a compact mass that points away from the lumen of the pharynx. In the horse, the auditory tube opens into the guttural pouch and the tonsils are diffuse and raised slightly. Horses are unable to breathe through the mouth as the free apex of the rostral epiglottis lies dorsal to the soft palate in a normal horse. In ruminants, the tonsils are a compact mass that points towards the lumen of the pharynx.
Pharyngeal arches
Pharyngeal arches are characteristic features of vertebrates whose origin can be traced back through chordates to basal deuterostomes who also share endodermal outpocketings of the pharyngeal apparatus. Similar patterns of gene expression can be detected in the developing pharynx of amphioxi and hemichordates. However, the vertebrate pharynx is unique in that it gives rise to endoskeletal support through the contribution of neural crest cells.
Pharyngeal jaws
Pharyngeal jaws are a "second set" of jaws contained within the pharynx of many species of fish, distinct from the primary (oral) jaws. Pharyngeal jaws have been studied in moray eels where their specific action is noted. When the moray bites prey, it first bites normally with its oral jaws, capturing the prey. Immediately thereafter, the pharyngeal jaws are brought forward and bite down on the prey to grip it; they then retract, pulling the prey down the eel's esophagus, allowing it to be swallowed.
Invertebrates
Invertebrates also have a pharynx. Invertebrates with a pharynx include the tardigrades, annelids and arthropods, and the priapulids (which have an eversible pharynx).
The "pharynx" of the nematode worm is a muscular food pump in the head, triangular in cross-section, that grinds food and transports it directly to the intestines. A one-way valve connects the pharynx to the excretory canal.
Additional images
| Biology and health sciences | Gastrointestinal tract | Biology |
2101485 | https://en.wikipedia.org/wiki/Beat%20%28acoustics%29 | Beat (acoustics) | In acoustics, a beat is an interference pattern between two sounds of slightly different frequencies, perceived as a periodic variation in volume whose rate is the difference of the two frequencies.
With tuning instruments that can produce sustained tones, beats can be readily recognized. Tuning two tones to a unison will present a peculiar effect: when the two tones are close in pitch but not identical, the difference in frequency generates the beating. The volume varies as in a tremolo as the sounds alternately interfere constructively and destructively. As the two tones gradually approach unison, the beating slows down and may become so slow as to be imperceptible. As the two tones get further apart, their beat frequency starts to approach the range of human pitch perception, the beating starts to sound like a note, and a combination tone is produced.
Mathematics and physics of beat tones
This phenomenon is best known in acoustics or music, though it can be found in any linear system:"According to the law of superposition, two tones sounding simultaneously are superimposed in a very simple way: one adds their amplitudes".If a graph is drawn to show the function corresponding to the total sound of two strings, it can be seen that maxima and minima are no longer constant (as when a pure note is played), but change over time: when the two waves are nearly 180 degrees out of phase the maxima of one wave cancel the minima of the other, whereas when they are nearly in phase their maxima sum up, raising the perceived volume.
It can be proven (with the help of a sum-to-product trigonometric identity) that the sum of two unit-amplitude sine waves can be expressed as a carrier wave of frequency whose amplitude is modulated by an envelope wave of frequency :
Because every other burst in the modulation pattern is inverted, each peak is replaced by a trough and vice versa. The envelope is perceived to have twice the frequency of the modulating cosine, which means the audible beat frequency (if it is in the audible range) is:
Monaural beats
"Monaural beats are when there is only one tone that pulses on and off in a specific pattern. With only one tone (as opposed to two tones with binaural beats), your brain has a much easier time adjusting and there is no need to balance separate tones.
Monaural beats are combined into one sound before they actually reach the human ear, as opposed to formulated in part by the brain itself, which occurs with a binaural beat.
This means that monaural beats can be used effectively via either headphones or speakers. It also means that those without two ears can listen to and receive the benefits." - Ebonie Allard
Binaural beats
"Binaural beats were first discovered by physicist Heinrich Wilhelm Dove in 1839. They are an auditory illusion that is perceived when two different pure-tone sine waves, both with frequencies lower than 1500 Hz and less than a 40 Hz difference between them, are presented to a listener dichotically (one through each ear). This means that they are best listened to through headphones." - Ebonie Allard
For example, if a 530 Hz pure tone is presented to a subject's right ear, while a 520 Hz pure tone is presented to the subject's left ear, the listener will hear beating at a rate of 10 Hz, just as if the two tones were presented monaurally, but the beating will have an element of lateral motion as well.
Binaural-beat perception originates in the inferior colliculus of the midbrain and the superior olivary complex of the brainstem, where auditory signals from each ear are integrated and precipitate electrical impulses along neural pathways through the reticular formation up the midbrain to the thalamus, auditory cortex, and other cortical regions.
According to a 2023 systematic review, studies have investigated some of the claimed positive effects in the areas of cognitive processing, affective states (like anxiety), mood, pain perception, meditation and relaxation, mind wandering, creativity, but the techniques were not comparable and results were inconclusive. Out of fourteen studies reviewed, five reported results in line with the brainwave entrainment hypothesis, eight studies reported contradictory, and one had mixed results. The authors recommend standardization in study approaches for future studies so results may be more effectively compared.
Uses
Musicians commonly use interference beats objectively to check tuning at the unison, perfect fifth, or other simple harmonic intervals. Piano and organ tuners use a method involving counting beats, aiming at a particular number for a specific interval.
The composer Alvin Lucier has written many pieces that feature interference beats as their main focus. Italian composer Giacinto Scelsi, whose style is grounded on microtonal oscillations of unisons, extensively explored the textural effects of interference beats, particularly in his late works such as the violin solos Xnoybis (1964) and L'âme ailée / L'âme ouverte (1973), which feature them prominently (Scelsi treated and notated each string of the instrument as a separate part, so that his violin solos are effectively quartets of one-strings, where different strings of the violin may be simultaneously playing the same note with microtonal shifts, so that the interference patterns are generated). Composer Phill Niblock's music is entirely based on beating caused by microtonal differences. Computer engineer Toso Pankovski invented a method based on auditory interference beating to screen participants in online auditory studies for headphones and dichotic context (whether the stereo channels are mixed or completely separated).
Amateur radio enthusiasts use the terms "zero-beating" or "zero-beat" for precisely tuning to a desired carrier wave frequency by manually reducing the number of interference beats, fundamentally the same tuning process used by musicians.
Sample
| Physical sciences | Waves | Physics |
2101553 | https://en.wikipedia.org/wiki/Saline%20water | Saline water | Saline water (more commonly known as salt water) is water that contains a high concentration of dissolved salts (mainly sodium chloride). On the United States Geological Survey (USGS) salinity scale, saline water is saltier than brackish water, but less salty than brine. The salt concentration is usually expressed in parts per thousand (permille, ‰) and parts per million (ppm). The USGS salinity scale defines three levels of saline water. The salt concentration in slightly saline water is 1,000 to 3,000 ppm (0.1–0.3%); in moderately saline water is 3,000 to 10,000 ppm (0.3–1%); and in highly saline water is 10,000 to 35,000 ppm (1–3.5%). Seawater has a salinity of roughly 35,000 ppm, equivalent to 35 grams of salt per one liter (or kilogram) of water. The saturation level is only nominally dependent on the temperature of the water. At one liter of water can dissolve about 357 grams of salt, a concentration of 26.3 percent by weight (% w/w). At (the boiling temperature of pure water), the amount of salt that can be dissolved in one liter of water increases to about 391 grams, a concentration of 28.1% w/w.
Properties
At , saturated sodium chloride brine is about 28% salt by weight. At , brine can only hold about 26% salt. At 20 °C one liter of water can dissolve about 357 grams of salt, a concentration of 26.3%.
The thermal conductivity of seawater (3.5% dissolved salt by weight) is 0.6 W/mK at . The thermal conductivity decreases with increasing salinity and increases with increasing temperature.
The salt content can be determined with a salinometer.
Density ρ of brine at various concentrations and temperatures from can be approximated with a linear equation:
where the values of an are:
Electrolysis
About four percent of hydrogen gas produced worldwide is created by electrolysis. The majority of this hydrogen produced through electrolysis is a side product in the production of chlorine.
2 NaCl(aq) + 2 H2O(l) → 2 NaOH(aq) + H2(g) + Cl2(g)
| Physical sciences | Water: General | Earth science |
2102802 | https://en.wikipedia.org/wiki/Rockslide | Rockslide | A rockslide is a type of landslide caused by rock failure in which part of the bedding plane of failure passes through compacted rock and material collapses en masse and not in individual blocks. Note that a rockslide is similar to an avalanche because they are both slides of debris that can bury a piece of land. While a landslide occurs when loose dirt or sediment falls down a slope, a rockslide occurs only when solid rocks are transported down slope. The rocks tumble downhill, loosening other rocks on their way and smashing everything in their path. Fast-flowing rock slides or debris slides behave similarly to snow avalanches, and are often referred to as rock avalanches or debris avalanches.
Definition
The term landslide refers to a variety of mass wasting events (geologic slope failures) that include slumps, slides, falls, and flows. The two major types of slides are rotational slides and translational slides. Rockslides are a type of translational event since the rock mass moves along a roughly planar surface with little rotation or backward tilting. Rock slides are the most dangerous form of mass-wasting because they incorporate a sudden, incredibly fast-paced release of bedrock along a uniform plane of weakness. These uniform weaknesses are key to identifying rock slides because unlike slumps, flows, or falls, the failed material moves in a fairly uniform direction over a layer of solid, pre-existing rock. Rock may break down while falling during rockslides.
The sudden, rapid release of material found in rock slides combined with the sheer size and weight of the material that is falling is what gives these events the potential to have devastating effects on human life and infrastructure. Rock slides are very common in the over steepened canyons and drainages of Idaho, particularly in those areas like the Salmon River Canyon where more than 5,000 feet of elevation may exist between the ridge tops and the canyon bottoms.
Causes
Mass-wasting occurs whenever the gravitational force on the rock exceeds the ability of the slope to resist that force. Therefore, anything that erodes or impedes the mountain's ability to resist this force may be one of the causes of mass-wasting. While a major event such as an earthquake can cause large rockslides to happen, a majority of slides occur due to a combination of gravitational pressure and erosional influences. Anthropogenic activities, such as altering the geometry of a hill through excavation, can also change the stress state, contributing to slope instability.
Amongst these erosional properties, water is arguably the most effective geologic agent that causes mass-wasting events to occur. Water aids in the downslope movement of surface material by adding weight to the soil and by filling pores which tends to push apart individual grains, decreasing the resistance of the material to movement. While these processes can cause a slide to happen, the speed and potential devastation of a rockslide is often determined by the severity of steepness presented by the failing slope.
Disaster prevention
With increasing populations in rural areas around the world, the hazards presented by potential rock slides are becoming more of a pressing issue moving forward. Luckily, individuals working in the fields of geologic science and engineering continue to perfect methods of rock slide detection, assessment, and warning. New earth observation tools have supplied a much enhanced ability to detect potential rock slide hazards. Analysis of sequential InSAR and LiDAR data provides a very valuable regional view of slope movement. Once susceptible areas are discovered, detailed analysis can be carried out at the specific site. These assessments are used to determine the amount of material that will be released as well as the speed at which this material will be transported.
Once a site is deemed hazardous, different types of geologic engineering techniques are used in order to prevent the compromised slope from failing. Some of these designs are listed below.
Wire meshing
A corrosion resistant material that is installed at the crest and foot of the slope. This insures that any falling debris is trapped behind the mesh.
Retaining walls
One of the oldest forms of ground engineering, retaining walls are built in order to neutralize the effects of unstable slopes by holding fallen rocks and soil back from roads and other structures.
Soil nailing
An economical method of constructing a form of retaining wall from the top of a slope down. In this process closely spaced steel tendons are drilled into the soil. These nails significantly increase the cohesion of soil through the ability of these rods to carry tensile loads. these steel beams are usually reinforced through the use of welded wire mesh.
Rock bolting
Rock bolts are always a primary means of reinforcement, Bolts are placed in a specific pattern in order to transfer the faces load from the slopes exterior, to its interior.
| Physical sciences | Geomorphology: General | Earth science |
2105550 | https://en.wikipedia.org/wiki/Bradymetabolism | Bradymetabolism | Bradymetabolism refers to organisms with a high active metabolism and a considerably slower resting metabolism. Bradymetabolic animals can often undergo dramatic changes in metabolic speed, according to food availability and temperature. Many bradymetabolic creatures in deserts and in areas that experience extreme winters are capable of "shutting down" their metabolisms to approach near-death states, until favorable conditions return
(see hibernation and estivation).
Several variants of bradymetabolism exists. In mammals, the animals normally have a fairly high metabolism, only dropping to low levels in times of little food. In most reptiles, the normal metabolic rate is quite low, but can be raised when needed, typically in short bursts of activity in connection with capturing prey.
Etymology
The term is from Greek brady (βραδύ) "slow" and metaballein (μεταβάλλειν) "turn quickly."
| Biology and health sciences | Basics | Biology |
30681089 | https://en.wikipedia.org/wiki/Mandible | Mandible | In jawed vertebrates, the mandible (from the Latin mandibula, 'for chewing'), lower jaw, or jawbone is a bone that makes up the lowerand typically more mobilecomponent of the mouth (the upper jaw being known as the maxilla).
The jawbone is the skull's only movable, posable bone, sharing joints with the cranium's temporal bones. The mandible hosts the lower teeth (their depth delineated by the alveolar process). Many muscles attach to the bone, which also hosts nerves (some connecting to the teeth) and blood vessels. Amongst other functions, the jawbone is essential for chewing food.
Owing to the Neolithic advent of agriculture (), human jaws evolved to be smaller. Although it is the strongest bone of the facial skeleton, the mandible tends to deform in old age; it is also subject to fracturing. Surgery allows for the removal of jawbone fragments (or its entirety) as well as regenerative methods. Additionally, the bone is of great forensic significance.
Structure
In humans, the mandible is the largest and lowest bone in the facial skeleton. It is the only movable bone of the skull (discounting the vibrating ossicles of the middle ear). It is connected to the skull's temporal bones by the temporomandibular joints. In addition to simply opening and closing, the jawbone can articulate side to side as well as forward and back.
Components
The mandible consists of:
The body, curving anteriorly like a horseshoe
Two rami (), rising from the posterior body, forming the (nearly square) angle of the mandible
Body
The body of the mandible is curved, and the front part gives structure to the chin. It has two surfaces and two borders. From the outside, the mandible is marked in the midline by a faint ridge, indicating the mandibular symphysis, the line of junction of the two halves of the mandible. This ridge divides below and encloses a triangular eminence, the mental protuberance (the chin), the base of which is depressed in the center but raised on both sides to form the mental tubercle. Just above this, on both sides, the mentalis muscles attach to a depression called the incisive foramen. Vertically midway on either side of the body, below the second premolar tooth, is the mental foramen, through which the mental nerve and blood vessels pass. Running backward and upward from each mental tubercle is a faint ridge, the oblique line, which is continuous with the anterior border of the ramus. Attached to this ridge is the masseter muscle (which covers most of the ramus and is a muscle of mastication), the depressor labii inferioris and depressor anguli oris (which support the mouth), and the platysma (extending down over much of the neck).
From the inside, the mandible appears concave. On either side of the lower symphysis is the mental spine (which can be faint or fused into one), to which the genioglossus (the inferior muscle of the tongue) attaches; the geniohyoid muscle attaches to the lower mental spine. Above the mental spine, a median foramen and furrow can line the symphysis. Below the mental spine is an oval depression (the digastric fossa of the mandible) where the digastric muscle attaches. Extending backward and upward on either side from the lower symphysis is a ridge called the mylohyoid line, where the mylohyoid muscle attaches; a small part of the superior pharyngeal constrictor muscle attaches to the posterior ridge, near the alveolar margin. Above the anterior ridge, the sublingual gland rests against a smooth triangular area, and below the posterior ridge, the submandibular gland rests in an oval depression.
Borders
The superior or alveolar border, wider behind than in front, is hollowed into cavities, for the reception of the teeth; these cavities are sixteen in number and vary in depth and size according to the teeth which they contain. To the outer lip of the superior border, on either side, the buccinator is attached as far forward as the first molar tooth.
The inferior border is rounded, longer than the superior, and thicker in front than behind; at the point where it joins the lower border of the ramus a shallow groove; for the facial artery, may be present. It is called the notch for the facial vessels or antegonial notch.
Ramus
The ramus of the human mandible has four sides, two surfaces, four borders, and two processes. On the outside, the ramus is flat and marked by oblique ridges at its lower part. It gives attachment throughout nearly the whole of its extent to the masseter muscle.
On the inside at the center there is an oblique mandibular foramen, for the entrance of the inferior alveolar vessels and nerve. The margin of this opening is irregular; it presents in front a prominent ridge, surmounted by a sharp spine, the lingula of the mandible, which gives attachment to the sphenomandibular ligament; at its lower and back part is a notch from which the mylohyoid groove runs obliquely downward and forward, and lodges the mylohyoid vessels and nerve. Behind this groove is a rough surface, for the insertion of the medial pterygoid muscle. The mandibular canal runs obliquely downward and forward in the ramus, and then horizontally forward in the body, where it is placed under the alveoli, with small openings for nerves. On arriving at the incisor teeth, it turns back to communicate with the mental foramen, giving off two small canals which run to the cavities containing the incisor teeth. In the posterior two-thirds of the bone the canal is situated nearer the internal surface of the mandible; and in the anterior third, nearer its external surface. It contains the inferior alveolar vessels and nerve, from which branches are distributed to the teeth.
Borders
The lower border of the ramus is thick, straight, and continuous with the inferior border of the body of the bone. At its junction with the posterior border is the angle of the mandible, which may be either inverted or everted and is marked by rough, oblique ridges on each side, for the attachment of the masseter laterally, and the medial pterygoid muscle medially; the stylomandibular ligament is attached to the angle between these muscles. The anterior border is thin above, thicker below, and continuous with the oblique line.
The region where the lower border meets the posterior border is the angle of the mandible.
The posterior border is thick, smooth, rounded, and covered by the parotid gland. The upper border is thin, and is surmounted by two processes, the coronoid in front and the condyloid behind, separated by a deep concavity, the mandibular notch.
Processes
The coronoid process is a thin, triangular eminence, which is flattened from side to side and varies in shape and size.
The condyloid process is thicker than the coronoid, and consists of two portions: the mandibular condyle, and the constricted portion which supports it, the neck. The condyle is the most superior part of the mandible and is part of the temporomandibular joint.
The mandibular notch, separating the two processes, is a deep semilunar depression and is crossed by the masseteric vessels and nerve.
Foramina
The mandible has two main holes (foramina), found on both its left and right sides:
The mandibular foramen, is above the mandibular angle in the middle of each ramus.
The mental foramen sits on either side of the mental protuberance (chin) on the body of mandible, usually inferior to the apices of the mandibular first and second premolars. As mandibular growth proceeds in young children, the mental foramen alters in direction of its opening from anterior to posterosuperior. The mental foramen allows the entrance of the mental nerve and blood vessels into the mandibular canal.
Nerves
The inferior alveolar nerve (IAN), a branch of the mandibular nerve (itself a major division of the cranium's trigeminal nerve), enters the mandibular foramen and runs forward in the mandibular canal, supplying sensation to the gums and teeth. Before passing through the mental foramen, the nerve divides into two terminal branches: incisive and mental nerves. The incisive nerve runs forward in the mandible and supplies the anterior teeth. The mental nerve exits the mental foramen and supplies sensation to the chin and lower lip.
Variation
Males generally have squarer, stronger, and larger mandibles than females. The mental protuberance is more pronounced in males but can be visualized and palpated in females.
Rarely, a bifid IAN may be present, resulting in a second and more inferiorly placed mandibular foramen. This can be detected by noting a doubled mandibular canal via radiograph.
Function
The mandible forms the lower jaw and holds the lower teeth in place. It articulates with the left and right temporal bones at the temporomandibular joints.
The condyloid process, the superior (upper) and posterior projection from the ramus, makes the temporomandibular joint with the temporal bone. The coronoid process, superior and anterior projection from the ramus. This provides attachment to the temporal muscle.
Teeth sit in the upper part of the body of the mandible. The frontmost part of teeth is more narrow and holds front teeth. The back part holds wider and flatter (albeit grooved) teeth primarily for chewing food.
The word mandible derives from the Latin word mandibula 'jawbone' (literally, 'used for chewing'), from mandere 'to chew' and -bula (instrumental suffix).
In addition to mastication, the joint of the jawbone enables actions such speech and yawning, while playing a more subtle role in activities such as kissing and breathing.
Phylogeny
The mandible of vertebrates evolved from Meckel's cartilage, left and right segments of cartilage which supported the anterior branchial arch in early fish. Fish jaws surface in species of the large arthrodire genus Dunkleosteus (), which crushed prey with their quickly articulating mouths. The lower jaw of cartilaginous fish, such as sharks, is composed of a cartilagenous structure homologous with Meckel's cartilage. This also remains a significant element of the jaw in some primitive bony fish, such as sturgeons. In reptiles, Meckel's cartilage ossifies into the (multiple) bones of the lower jaw, while mammals of the Cretaceous (145–66 Mya) had both Meckel's cartilage and a mandible.
In lobe-finned fishes and the early fossil tetrapods, the bone homologous to the mandible of mammals is merely the largest of several bones in the lower jaw. In such animals, it is referred to as the dentary bone or os dentale, and forms the body of the outer surface of the jaw. It is bordered below by a number of splenial bones, while the angle of the jaw is formed by a lower angular bone and a suprangular bone just above it. The inner surface of the jaw is lined by a prearticular bone, while the articular bone forms the articulation with the skull proper. A set of three narrow coronoid bones lie above the prearticular bone. As the name implies, the majority of the teeth are attached to the dentary, but there are commonly also teeth on the coronoid bones, and sometimes on the prearticular.
Most vertebrates exhibit a simpler scheme, as bones have either fused or vanished. In teleosts, only the dentary, articular, and angular bones remain, while in living amphibians, the dentary is accompanied only by the prearticular, and, in salamanders, one of the coronoids. The lower jaw of reptiles has only a single coronoid and splenial, but retains all the other primitive bones except the prearticular and the periosteum. In birds, the various bones have fused into a single structure. In mammals, most have disappeared, leaving only the mandible. As a result, there is only articulation between the mandible and temporal bones, as opposed to articulation between articular and quadrate bones. An intermediate stage can be seen in some therapsids, in which both points of articulation are present. Aside from the dentary, only few other bones of the lower jaw remain in mammals; the former articular and quadrate bones survive as the malleus and the incus of the middle ear.
In recent human evolution, both the oral cavity and jaws have shrunk in correspondence with the Neolithic-era shift from hunter-gatherer lifestyles towards agriculture and settlement, dated to . This has led to orthodontic malocclusions.
Development
The mandible forms as a bone (ossifies) from Meckel's cartilage, which forms the cartilaginous bar of the mandibular arch and, dorsally, parts of the middle ear. The two sides of the jawbone are inferiorly fused at the mandibular symphysis (the chin) during the first year of life. The cartilage of the ramus is replaced by fibrous tissue, which persists to form the sphenomandibular ligament. Between the lingula and the canine tooth the cartilage disappears, while the portion of it below and behind the incisor teeth becomes ossified and incorporated with this part of the mandible.
About the sixth week of fetal life, intramembranous ossification takes place in the membrane covering the outer surface of the ventral end of Meckel's cartilage, and each half of the bone is formed from a single center which appears, near the mental foramen. By the tenth week, the portion of Meckel's cartilage which lies below and behind the incisor teeth is surrounded and invaded by the dermal bone (also known as the membrane bone). Somewhat later, accessory nuclei of cartilage make their appearance, as
a wedge-shaped nucleus in the condyloid process and extending downward through the ramus;
a small strip along the anterior border of the coronoid process;
smaller nuclei in the front part of both alveolar walls and along the front of the lower border of the bone.
These accessory nuclei possess no separate ossific centers but are invaded by the surrounding dermal bone and undergo absorption. The inner alveolar border, usually described as arising from a separate ossific center (splenial center), is formed in the human mandible by an ingrowth from the main mass of the bone.
Aging
At birth, the body of the bone is a mere shell, containing the sockets of the two incisor, the canine, and the two deciduous molar teeth, imperfectly partitioned off from one another. The mandibular canal is of large size and runs near the lower border of the bone; the mental foramen opens beneath the socket of the first deciduous molar tooth. The angle is obtuse (175°), and the condyloid portion is nearly in line with the body. The coronoid process is of comparatively large size, and projects above the level of the condyle.
After birth, the two segments of the bone become joined at the symphysis, from below upward, in the first year; but a trace of separation may be visible in the beginning of the second year, near the alveolar margin. The body becomes elongated in its whole length, but more especially behind the mental foramen, to provide space for the three additional teeth developed in this part. The depth of the body increases owing to increased growth of the alveolar part, to afford room for the roots of the teeth, and by thickening of the subdental portion which enables the jaw to withstand the powerful action of the masticatory muscles; but, the alveolar portion is the deeper of the two, and, consequently, the chief part of the body lies above the oblique line. The mandibular canal, after the second dentition, is situated just above the level of the mylohyoid line; and the mental foramen occupies the position usual to it in the adult. The angle becomes less obtuse, owing to the separation of the jaws by the teeth; about the fourth year it is 140°. The fibrocartilage of the mandibular symphysis fuses together in early childhood.
In the adult, the alveolar and subdental portions of the body are usually of equal depth. The mental foramen opens midway between the upper and lower borders of the bone, and the mandibular canal runs nearly parallel with the mylohyoid line. The ramus is almost vertical in direction, the angle measuring from 110° to 120°, also the adult condyle is higher than the coronoid process and the sigmoid notch becomes deeper. The adult mandible is the skull's largest and strongest bone.
In old age, the bone can become greatly reduced in volume where there is a loss of teeth, and consequent resorption of the alveolar process and interalveolar septa. Consequently, the chief part of the bone is below the oblique line. The mandibular canal, with the mental foramen opening from it, is closer to the alveolar border. The ramus is oblique in direction, the angle measures about 140°, and the neck of the condyle is more or less bent backward.
Clinical significance
The posterior of the mandible is notoriously resistant to the full effects of local anesthesia. The IAN provides sensory innervation to much of the mandible and its teeth, making it a target of block anesthesia. Injecting the nerve is challenging due to the amount of surrounding soft tissue. American surgeon William Stewart Halsted developed a technique using a syringe and cocaine which was performed successfully by 1885.
Fracture
One fifth of facial injuries involve a mandibular fracture. Mandibular fractures are often accompanied by a 'twin fracture' on the opposite side. There is no universally accepted treatment protocol, as there is no consensus on the choice of techniques in a particular anatomical shape of mandibular fracture clinic. A common treatment involves attachment of metal plates to the fracture to assist in healing.
Dislocation and displacement
The mandible may be dislocated anteriorly (to the front) and inferiorly (downwards) but very rarely posteriorly (backwards). The articular disk of the temporomandibular joint prevents the mandible from moving posteriorly, making the condylar neck particularly vulnerable to fractures. Further, various jawbone damage can cause temporomandibular joint dysfunction, with symptoms including pain and inflammation.
The jawbone can also become deviated in mandibular lateral displacement, a condition which can offset facial symmetry and cause posterior crossbite.
Resorption
The mandibular alveolar process can become resorbed when completely edentulous in the mandibular arch (occasionally noted also in partially edentulous cases). This resorption can occur to such an extent that the mental foramen is virtually on the superior border of the mandible, instead of opening on the anterior surface, changing its relative position. However, the more inferior body of the mandible is not affected and remains thick and rounded. With age and tooth loss, the alveolar process is absorbed so that the mandibular canal becomes nearer the superior border. Sometimes with excessive alveolar process absorption, the mandibular canal disappears entirely and deprives the IAN of its bony protection, although soft tissue continues to guard the nerve.
Mandibulectomy
The surgical removal (resection) of all or part of the jawbone is known as a mandibulectomy. The removal of a small portion is known as partial mandibulectomy and a larger portion segmental mandibulectomy. This can be performed in response to cancer (i.e. tumor removal), infection, injury, or osteonecrosis. The removed portion can be replaced with metal plating or bone from elsewhere in the body. Oral muscles tend to work differently after the procedure, requiring therapy to relearn operations such as eating and speaking. During recovery, a feeding tube is utilized, and sometimes a tracheotomy is also performed to maintain respiration in case of swollen muscles. In a technique illustrated in the 19th century, the flesh of the face is incised along the inferior of the mandible and peeled upward for the bone's removal.
Complications can involve difficulties with free flap transfer and airway management. Additional side effects include pain, infection, numbness, and (rarely, fatal) bleeding. Even successful surgeries can result in deformity, with an extreme version being referred to as the Andy Gump deformity after the comic book character, whose design apparently lacks a jaw; proposed reconstruction methods include implanting synthetic material, potentially involving 3D printing.
Regeneration
Bone loss (as in osteoporosis) can be mitigated in the jawbone via bone grafting, which is sometimes performed to support dental implants (replacing teeth individually or in groups).
Mandibular prosthetics date back to ancient Egypt and China, but significant advancements were made in the late 19th century with new techniques for attaching prosthetics to a depreciated jawbone as well as bone grafting.
In 2010, the first successful face transplant was conducted on a Spanish farmer after a self-inflicted gun accident; this included the replacement of the entire mandible.
Forensic medicine
The mandible can provide forensic evidence because its form changes over a person's life, and this can be used to determine a deceased person's age.
Dental remains of Nazi leader Adolf Hitler, including part of a mandible with teeth, were the solitary physical evidence used to confirm his death in 1945.
In culture
In the Hebrew Bible and Christian Old Testament Book of Judges, Samson used a donkey's jawbone to kill a thousand Philistines.
As early as 1900, the phrase jaw-dropping was used as an adjective to describe a condition of shock in humans, e.g. when someone's mouth suddenly hangs agape in response to something. The exaggerated visual gag of a jaw dropping to the floor was a trademark of American animation director Tex Avery, who would often employ it when the Big Bad Wolf spies a sexually attractive woman.
Gobstoppers, a type of hard candy, are known in North America as jawbreakers due to the fracturing risk they impose on teeth.
Owing in part to the forensic evidence of Hitler's death being limited to his dental remains (including a jawbone fragment broken and burnt around the alveolar process), some fringe accounts (bolstered by the Soviet Union, which captured Berlin in 1945) allege that Hitler faked his death (ostensibly along with Eva Braun).
In later decades, American real-estate businessman Fred Trump had a partial mandibulectomy which caused a conspicuous deformity. In his fight against cancer, American film critic Roger Ebert had a partial mandibulectomy in 2006, in addition to later surgeries.
Additional images
| Biology and health sciences | Human anatomy | Health |
30683135 | https://en.wikipedia.org/wiki/Yellow-bellied%20sea%20snake | Yellow-bellied sea snake | The yellow-bellied sea snake (Hydrophis platurus) is an extremely venomous species of snake from the subfamily Hydrophiinae (the sea snakes) found in tropical oceanic waters around the world except for the Atlantic Ocean. For many years, it was placed in the monotypic genus Pelamis, but recent molecular evidence indicates it lies within the genus Hydrophis.
Taxonomy
In 1766, Linnaeus published the original description of the yellow-bellied sea snake, naming it Anguis platura (Anguis meaning snake). In 1803, François Marie Daudin created the new genus Pelamis and assigned this species to it, referring to it as Pelamis platuros. In 1842, Gray described what he thought was a new species and called it Pelamis ornata (subsequently P. ornata became a synonym of P. platura). The commonly used genus name Pelamis is derived from the Ancient Greek word for "tunny fish", which presumably refers to the habitat or what Daudin thought they ate. The specific name platurus is a combination of the Ancient Greek words platys "flat" and oura "tail", referring to the flattened tail. The word Pelamis is a feminine noun and means young or small tunny fish. In 1872, Stoliczka introduced the name Pelamis platurus (still the most used scientific name by scientists today), but used the incorrect ending -us instead of -a which a feminine noun requires. A few recent examples exist of scientists' beginning to use the grammatically correct name Pelamis platura, e.g., Bohme 2003 and the Reptile Database with its page headed Pelamis platura (Linnaeus, 1766), which includes an extensive synonymy of the different scientific names which have been used for the yellow-bellied sea snake. The same rules apply for the most recent taxonomic name of Hydrophis platurus.
To further complicate the nomenclature, the taxonomic status of sea snakes is still under review, with recent authors suggesting a dismantling of monotypic genera, such as Pelamis, in favour of a single genus, Hydrophis, in order to reduce paraphyly and better reflect phylogenetic relationships.
Other common names are yellowbelly sea snake or pelagic sea snake.
Evolution
Sea snakes are a monophyletic group (Hydrophiinae) that diverged from the front-fanged Australasian venomous snakes (Elapidae) about 10 million years ago. The yellow-bellied sea snake is a part of the rapidly radiating Hydrophis group.
Description
The yellow-bellied sea snake, as the name implies, has a distinctive bicolor pattern with a yellow underbelly and brown back, making it easily distinguishable from other sea snake species. Yellow-bellied sea snakes, like many other species of sea snake, are fully adapted to living their whole lives at sea: mating, eating and giving birth to live young (ovoviviparous). Adaptations to aquatic life include the reduced ventral scale size, laterally compressed body and paddle-tail for swimming, valved nostrils and palatine seal for excluding seawater, and cutaneous gas exchange for prolonging dive times. This species can uptake up to 33% of its oxygen requirements through the skin while diving and swimming at the surface of the water. Sea snakes also have a special salt gland located in the lower jaw that was formerly believed to filter out salt from the surrounding seawater but has been found not to be used for that purpose, as sea snakes drink fresh water only.
See snake scales for terminology used here
Morphology
The body of this snake is compressed, with the posterior less than half the diameter of the neck; the body scales are juxtaposed, subquadrangular in shape, and in 23–47 rows around the thickest part of the body; ventral scales, 264–406 in number, are very small and, if distinct, divided by a latitudinal
groove, but usually are indistinguishable from adjacent body scales. The head is narrow, with an elongated snout; head shields are entire, nostrils are superior, and nasal shields are in contact with one another; the prefrontal scale is in contact with second upper labial; one or two preoculars, two or three postoculars, and two or three small anterior temporals are present; seven or eight upper labials are found, with four or five below the eye, but separated from the border by a subocular. Colors of the snake are variable, but most often distinctly bicolored, black above, yellow or brown below, with the dorsal and ventral colors sharply demarcated from one another; ventrally, there may be a series of black spots or bars on the yellow or brown background, or the yellow may extend dorsally so there is only a narrow middorsal black stripe, or a series of black crossbars. Total length for males is up to , for females up to ; tail length for males is up to , females up to .
Distribution and habitat
The yellow-bellied sea snake is one of the most widely distributed snakes in the world. It is completely pelagic and is often observed on oceanic drift lines, using surface currents and storms to move around the ocean. Their distribution appears to be largely determined by favourable water temperatures, oceanic currents and recent formation of land bridges that have blocked farther dispersal.
The yellow-bellied sea snake has an extensive distribution covering the entire tropical Indo-Pacific, as well as extending to Costa Rica, southern California, and northern Peru. It is the only sea snake to have reached the Hawaiian Islands. The favoured habitat for hunting and reproduction includes free floating mats of sea kelp occurring in the Indian Ocean. The species is the most commonly beached sea snake on the coast of Southwest Australia, including records at beaches near metropolitan areas. It is also reported from Christmas Island and Cocos (Keeling) Islands (Australia).
The yellow-bellied sea snake requires a minimum of to survive, long-term. However, the species has been reported in colder waters of the Pacific, such as the coasts of southern California, Mexico, Tasmania, and New Zealand, the latter being a country that would otherwise be free of snakes were it not for the infrequent strandings of yellow-bellied sea snakes and banded sea kraits. Nonetheless, these wayward individuals make the yellow-bellied sea snake the most commonly-seen snake (and sea snake) in New Zealand, to the degree that the species is considered native (indigenous) to New Zealand and worthy of protection under the Wildlife Act 1953. These colder water occurrences are believed to be linked to El Niño, among other severe weather events, possibly creating unusually strong, new ocean currents that transport the snakes far off-course. In October 2015, beached yellow-bellied sea snakes were reported and photographed on beaches in Ventura County, California, well outside their normal range, for the first time in 30 years. A few months later, in January 2016, a stranded individual was found in Coronado, California, washed-up on Coronado Beach's north end (better known as Dog Beach), just south of Naval Air Station North Island (NASNI). The specimen was subsequently transported to and examined at the Scripps Institute of Oceanography, La Jolla.
The yellow-bellied sea snake is the only sea snake to have been found in the Atlantic Ocean, although only in limited circumstances. The yellow-bellied sea snake's occurrence into the Atlantic is not considered a part of its native range, but rather a dispersal from its native Pacific range.
The yellow-bellied sea snake has been found in all the countries of Africa's eastern coast and all eastern islands, like Djibouti, Eritrea, Kenya, Madagascar, Mauritius, Mayotte, Mozambique, Réunion, Seychelles, Somalia, South Africa and Tanzania. On the African Atlantic coast they have been reported to occur in the Benguela Current, with specimens found along the coasts of South Africa and Namibia.
The yellow-bellied sea snake has also been found in the Colombian Caribbean four separate times, making it the only sea snake to be found in the Caribbean Sea. However these occurrences are believed to be the result of human activity, be it ship discharge, intentional release or via the Panama Canal, as it is not considered a part of their native range. This is due to the land bridge between North and South America, the Isthmus of Panama, which formed from about 10 million years ago to 3 million years ago (i.e., continental drift), acting as a dispersal barrier and preventing entry into the Caribbean Sea from the Pacific Ocean. The man-made Panama Canal has not made a crossing of the isthmus possible presumably because it is fresh water.
Due to the wide distribution of the species and relative lack of dispersal barriers, it has been assumed that individuals from different localities represent a single breeding population (i.e., high gene flow). However, a study that used haplotype networks in two populations from Costa Rica suggests that shallow genetic population structure exists, which reflects variation in colour patterns (brown and yellow in Golfo de Papagayo and completely yellow in Golfo Dulce).
Behaviour
Contrary to past beliefs, sea snakes require fresh water to survive and the yellow-bellied sea snake drinks precipitation that forms on the surface of sea water. This species has been reported to survive severe dehydration of up to 7 months during seasonal drought.
Yellow-bellied sea snakes breed in warm waters; they are ovoviviparous with a gestation period around 6 months. According to Ditmars, females bear live young in tidal pools. They move poorly on land due to their smaller belly scales that form a ventral keel. They are sometimes observed in large aggregations of thousands on the surface of the water in oceanic drift lines, which has been proposed as a strategy to catch prey. They hunt by floating on the surface of the water to attract pelagic fish that are seeking shelter; prey are captured via a backwards swimming motion and rapid lunge of the jaws. The ability to swim backwards is an unusual and distinguishing characteristic of this species. Heatwole proposed that these snakes find their prey by sensing the vibration generated by fish movement. These snakes are not considered aggressive but will defend themselves if threatened.
Venom
The venom of this species is highly potent, like that of other sea snakes. Bites are rare and the most common victims are fishermen who try to get them out of fishing nets. The subcutaneous of the venom is 0.067 mg/kg and the venom yield per bite is 1.0–4.0 mg. Yellow-bellied sea snake venom contains several different neurotoxins and two other isotoxins.
Antivenom
Sea snake venom can cause damage to skeletal muscle with consequent myoglobinuria, neuromuscular paralysis or direct renal damage. The venoms of significant species of sea snake are neutralised with Commonwealth Serum Laboratories Ltd (of Melbourne, Australia) sea snake (Enhydrina schistosa) antivenom. If that preparation is not available, tiger snake or polyvalent antivenom should be used. No deaths have been recorded from bites in Australian waters. The E. schistosa antivenom was tested specifically on Pelamus platurus, and it effectively neutralised the venom.
Cited references
Other sources
Ditmars, R.L. 1936. The Reptiles of North America. Doubleday, Doran & Co. New York. 476 pp.
Kropach, C. 1975 The yellow-bellied sea snake, Pelamis, in the eastern Pacific. pp. 185–213 in: Dunson, W., ed., The Biology of Sea Snakes. Univ. Park Press, Baltimore, xi + 530 pp.
Smith, M.A. 1943. The Fauna of British India, Ceylon and Burma, including the Whole of the Indo-Chinese Sub-region. Reptiles and Amphibians. Vol. III. – Serpentes. Taylor & Francis. London. 583 pp.
External links
Pelagic Sea Snakes and the animals that live on them, Life is Short but Snakes are Long
Reptiles described in 1766
Taxa named by Carl Linnaeus
Hydrophis
Reptiles of Pakistan
Reptiles of New Zealand
Reptiles of Western Australia
Snakes of Japan
Reptiles of Mexico
Reptiles of Iran
Reptiles of Guatemala
Reptiles of Hawaii
Snakes of Australia | Biology and health sciences | Snakes | Animals |
30688796 | https://en.wikipedia.org/wiki/Quercus%20oblongifolia | Quercus oblongifolia | Quercus oblongifolia, commonly known as the Mexican blue oak, Arizona blue oak, Blue live oak or Sonoran blue oak, is an evergreen small tree or large shrub in the white oak group.
Distribution
Quercus oblongifolia grows in high grasslands, canyons and mesas in southwestern United States (Texas, Arizona and New Mexico) and northwestern Mexico (Baja California Sur, Chihuahua, Coahuila, Sinaloa and Sonora states). Mexican blue oak is closely related to Engelmann oak "Quercus engelmannii" in Southern California. The two species may be conspecific.
Description
The Mexican blue oak is a small evergreen tree growing 5–8 metres (16–27 feet) tall with a rounded crown. At higher elevations it is typically a large shrub. The trunk is up to in diameter and the bark is light gray and densely furrowed. The twigs are yellowish brown and hairless with reddish brown buds. The leaves are small, alternate and oblong, with entire margins, leathery, bluish-green above and mid green below. The flowers appear in spring at the same time as the old leaves are being shed and new leaf growth starts. The male flowers form yellowish-green catkins and the female flowers are solitary or paired and grow in the leaf axils. The light brown acorns are ovoid or oblong, about long and lodged in scaly, bowl-shaped cups about one third the length of the nut.
Habitat
The Mexican blue oak is common at elevations of 1,200 to 1,800 m (4,000–6,000 ft). It is often found on thin sandy soils in semi-arid regions and is the dominant species in lower open oak woodland where it grows in association with Arizona white oak (Quercus arizonica) and Emory oak (Quercus emoryi). It is an important constituent of pinyon–juniper communities. where it grows in association with species of pine and juniper, Arizona rosewood (Vauquelinia californica), shrubby buckwheat (Eriogonum wrightii), catclaw mimosa (Mimosa aculeaticarpa), bullgrass (Muhlenbergia emersleyi), plains lovegrass (Eragrostis intermedia), fendlerbush (Fendlera rupicola) and wolftail (Lycurus phleoides).
Cultivation
Mexican blue oak "Quercus oblongifolia" popularity in landscaping has been increasing in California with many plantings at Apple Park and other Bay area plantings. The fast growth and beautiful blue foliage makes selections from trees in Arizona the best for cultivation.
| Biology and health sciences | Fagales | Plants |
3919387 | https://en.wikipedia.org/wiki/Remineralisation | Remineralisation | In biogeochemistry, remineralisation (or remineralization) refers to the breakdown or transformation of organic matter (those molecules derived from a biological source) into its simplest inorganic forms. These transformations form a crucial link within ecosystems as they are responsible for liberating the energy stored in organic molecules and recycling matter within the system to be reused as nutrients by other organisms.
Remineralisation is normally viewed as it relates to the cycling of the major biologically important elements such as carbon, nitrogen and phosphorus. While crucial to all ecosystems, the process receives special consideration in aquatic settings, where it forms a significant link in the biogeochemical dynamics and cycling of aquatic ecosystems.
Role in biogeochemistry
The term "remineralization" is used in several contexts across different disciplines. The term is most commonly used in the medicinal and physiological fields, where it describes the development or redevelopment of mineralized structures in organisms such as teeth or bone. In the field of biogeochemistry, however, remineralization is used to describe a link in the chain of elemental cycling within a specific ecosystem. In particular, remineralization represents the point where organic material constructed by living organisms is broken down into basal inorganic components that are not obviously identifiable as having come from an organic source. This differs from the process of decomposition which is a more general descriptor of larger structures degrading to smaller structures.
Biogeochemists study this process across all ecosystems for a variety of reasons. This is done primarily to investigate the flow of material and energy in a given system, which is key to understanding the productivity of that ecosystem along with how it recycles material versus how much is entering the system. Understanding the rates and dynamics of organic matter remineralization in a given system can help in determining how or why some ecosystems might be more productive than others.
Remineralization reactions
While it is important to note that the process of remineralization is a series of complex biochemical pathways [within microbes], it can often be simplified as a series of one-step processes for ecosystem-level models and calculations. A generic form of these reactions is shown by:
{Organic \ Matter} + Oxidant -> {Liberated\ Simple\ Nutrients} + \underset{Carbon~Dioxide}{CO2} + \underset{Water}{H2O}
The above generic equation starts with two reactants: some piece of organic matter (composed of organic carbon) and an oxidant. Most organic carbon exists in a reduced form which is then oxidized by the oxidant (such as ) into and energy that can be harnessed by the organism. This process generally produces , water and a collection of simple nutrients like nitrate or phosphate that can then be taken up by other organisms. The above general form, when considering as the oxidant, is the equation for respiration. In this context specifically, the above equation represents bacterial respiration though the reactants and products are essentially analogous to the short-hand equations used for multi-cellular respiration.
Electron acceptor cascade
The degradation of organic matter through respiration in the modern ocean is facilitated by different electron acceptors, their favorability based on Gibbs free energy law, and the laws of thermodynamics. This redox chemistry is the basis for life in deep sea sediments and determines the obtainability of energy to organisms that live there. From the water interface moving toward deeper sediments, the order of these acceptors is oxygen, nitrate, manganese, iron, and sulfate. The zonation of these favored acceptors can be seen in Figure 1. Moving downwards from the surface through the zonation of these deep ocean sediments, acceptors are used and depleted. Once depleted the next acceptor of lower favorability takes its place. Thermodynamically, oxygen represents the most favorable electron accepted but is quickly used up in the water sediment interface and concentrations extends only millimeters to centimeters down into the sediment in most locations of the deep sea. This favorability indicates an organism's ability to obtain higher energy from the reaction which helps them compete with other organisms. In the absence of these acceptors, organic matter can also be degraded through methanogenesis, but the net oxidation of this organic matter is not fully represented by this process. Each pathway and the stoichiometry of its reaction are listed in table 1.
Due to this quick depletion of in the surface sediments, a majority of microbes use anaerobic pathways to metabolize other oxides such as manganese, iron, and sulfate. It is also important to figure in bioturbation and the constant mixing of this material which can change the relative importance of each respiration pathway. For the microbial perspective please reference the electron transport chain.
Remineralisation in sediments
Reactions
A quarter of all organic material that exits the photic zone makes it to the seafloor without being remineralised and 90% of that remaining material is remineralised in sediments itself. Once in the sediment, organic remineralisation may occur through a variety of reactions. The following reactions are the primary ways in which organic matter is remineralised, in them general organic matter (OM) is often represented by the shorthand: .
Aerobic respiration
Aerobic respiration is the most preferred remineralisation reaction due to its high energy yield. Although oxygen is quickly depleted in the sediments and is generally exhausted centimeters from the sediment-water interface.
Anaerobic respiration
In instances in which the environment is suboxic or anoxic, organisms will prefer to utilize denitrification to remineralise organic matter as it provides the second largest amount of energy. In depths below where denitrification is favored, reactions such as Manganese Reduction, Iron Reduction, Sulfate Reduction, Methane Reduction (also known as Methanogenesis), become favored respectively. This favorability is governed by Gibbs Free Energy (ΔG). In a water body, sediment seabed, or soil, the sorting of these chemical reactions with depth in order of energy provided is called a redox gradient.
Redox zonation
Redox zonation refers to how the processes that transfer terminal electrons as a result of organic matter degradation vary depending on time and space. Certain reactions will be favored over others due to their energy yield as detailed in the energy acceptor cascade detailed above. In oxic conditions, in which oxygen is readily available, aerobic respiration will be favored due to its high energy yield. Once the use of oxygen through respiration exceeds the input of oxygen due to bioturbation and diffusion, the environment will become anoxic and organic matter will be broken down via other means, such as denitrification and manganese reduction.
Remineralisation in the open ocean
In most open ocean ecosystems only a small fraction of organic matter reaches the seafloor. Biological activity in the photic zone of most water bodies tends to recycle material so well that only a small fraction of organic matter ever sinks out of that top photosynthetic layer. Remineralisation within this top layer occurs rapidly and due to the higher concentrations of organisms and the availability of light, those remineralised nutrients are often taken up by autotrophs just as rapidly as they are released.
What fraction does escape varies depending on the location of interest. For example, in the North Sea, values of carbon deposition are ~1% of primary production while that value is <0.5% in the open oceans on average. Therefore, most of nutrients remain in the water column, recycled by the biota. Heterotrophic organisms will utilize the materials produced by the autotrophic (and chemotrophic) organisms and via respiration will remineralise the compounds from the organic form back to inorganic, making them available for primary producers again.
For most areas of the ocean, the highest rates of carbon remineralisation occur at depths between in the water column, decreasing down to about 1,200 m where remineralisation rates remain pretty constant at 0.1 μmol kg−1 yr−1. As a result of this, the pool of remineralised carbon (which generally takes the form of carbon dioxide) tends to increase in the photic zone.
Most remineralisation is done with dissolved organic carbon (DOC). Studies have shown that it is larger sinking particles that transport matter down to the sea floor while suspended particles and dissolved organics are mostly consumed by remineralisation. This happens in part due to the fact that organisms must typically ingest nutrients smaller than they are, often by orders of magnitude. With the microbial community making up 90% of marine biomass, it is particles smaller than the microbes (on the order of ) that will be taken up for remineralisation.
| Physical sciences | Geochemistry | Earth science |
3925077 | https://en.wikipedia.org/wiki/Binary%20pulsar | Binary pulsar | A binary pulsar is a pulsar with a binary companion, often a white dwarf or neutron star. (In at least one case, the double pulsar PSR J0737-3039, the companion neutron star is another pulsar as well.) Binary pulsars are one of the few objects which allow physicists to test general relativity because of the strong gravitational fields in their vicinities. Although the binary companion to the pulsar is usually difficult or impossible to observe directly, its presence can be deduced from the timing of the pulses from the pulsar itself, which can be measured with extraordinary accuracy by radio telescopes.
History
The binary pulsar PSR B1913+16 (or the "Hulse-Taylor binary pulsar") was first discovered in 1974 at Arecibo by Joseph Hooton Taylor, Jr. and Russell Hulse, for which they won the 1993 Nobel Prize in Physics. While Hulse was observing the newly discovered pulsar PSR B1913+16, he noticed that the rate at which it pulsed varied regularly. It was concluded that the pulsar was orbiting another star very closely at a high velocity, and that the pulse period was varying due to the Doppler effect: As the pulsar was moving towards Earth, the pulses would be more frequent; and conversely, as it moved away from Earth fewer would be detected in a given time period. One can think of the pulses like the ticks of a clock; changes in the ticking are indications of changes in the pulsars speed toward and away from Earth. Hulse and Taylor also determined that the stars were approximately equally massive by observing these pulse fluctuations, which led them to believe the other object was also a neutron star. Pulses from this system are now tracked to within 15 μs. (Note: Cen X-3 was actually the first "binary pulsar" discovered in 1971, followed by Her X-1 in 1972.)
The study of the PSR B1913+16 binary pulsar also led to the first accurate determination of neutron star masses, using relativistic timing effects. When the two bodies are in close proximity, the gravitational field is stronger, the passage of time is slowed – and the time between pulses (or ticks) is lengthened. Then as the pulsar clock travels more slowly through the weakest part of the field it regains time. A special relativistic effect, time dilation, acts around the orbit in a similar fashion. This relativistic time delay is the difference between what one would expect to see if the pulsar were moving at a constant distance and speed around its companion in a circular orbit, and what is actually observed.
Prior to the first observation of gravitational waves in 2015 and the operation of Advanced LIGO, binary pulsars were the only tools scientists had to detect evidence of gravitational waves; Einstein's theory of general relativity predicts that two neutron stars would emit gravitational waves as they orbit a common center of mass, which would carry away orbital energy and cause the two stars to draw closer together and shorten their orbital period. A 10-parameter model incorporating information about the pulsar timing, the Keplerian orbits and three post-Keplerian corrections (the rate of periastron advance, a factor for gravitational redshift and time dilation, and a rate of change of the orbital period from gravitational radiation emission) is sufficient to completely model the binary pulsar timing.
The measurements made of the orbital decay of the PSR B1913+16 system were a near perfect match to Einstein's equations. Relativity predicts that over time a binary system's orbital energy will be converted to gravitational radiation. Data collected by Taylor and Joel M. Weisberg and their colleagues of the orbital period of PSR B1913+16 supported this relativistic prediction; they reported in 1982 and subsequently that there was a difference in the observed minimum separation of the two pulsars compared to that expected if the orbital separation had remained constant. In the decade following its discovery, the system's orbital period had decreased by about 76 millionths of a second per year, indicating that the pulsar was approaching its maximum separation more than a second earlier than it would have if the orbit had remained the same. Subsequent observations continue to show this decrease.
Intermediate mass binary pulsar
An (IMBP) is a pulsar-white dwarf binary system with a relatively long spin period of around 10–200 ms consisting of a white dwarf with a relatively high mass of approximately
The spin periods, magnetic field strengths, and orbital eccentricities of IMBPs are significantly larger than those of low mass binary pulsars (LMBPs). As of 2014, there are fewer than 20 known IMBPs. Examples of IMBPs include PSR J1802−2124 and PSR J2222−0137.
The binary system PSR J2222−0137 has an orbital period of about 2.45 days and is found at a distance of 267 pc (approximately 870 light-years), making it the second closest known binary pulsar systems (as of 2014) and one of the closest pulsars and neutron stars. The relatively high-mass pulsar (1.831 0.010 ) has a companion star PSR J2222−0137 B with a minimum mass of approximately 1.3 solar masses (1.319 0.004 ). This meant the companion is a massive white dwarf (only about 8% of white dwarfs have a mass ), which would make the system an IMBP. Although initial measurements gave a mass of about 1 solar mass for the PSR J2222−0137 B, later observations showed that it is actually a high-mass white dwarf and also one of the coolest known white dwarfs, with a temperature less than 3,000 K.
PSR J2222−0137 B is likely crystallized, leading to this Earth-sized white dwarf being described as a "diamond-star", similar to the white dwarf companion of PSR J1719-1438, which lies about 4,000 light-years away.
Effects
Sometimes the relatively normal companion star of a binary pulsar will swell up to the point that it dumps its outer layers onto the pulsar. This interaction can heat the gas being exchanged between the bodies and produce X-ray light which can appear to pulsate, in a process called the X-ray binary stage. The flow of matter from one stellar body to another often leads to the creation of an accretion disk about the recipient star.
Pulsars also create a "wind" of relativistically outflowing particles, which in the case of binary pulsars can blow away the magnetosphere of their companions and have a dramatic effect on the pulse emission.
| Physical sciences | Stellar astronomy | Astronomy |
3925533 | https://en.wikipedia.org/wiki/Enumerative%20combinatorics | Enumerative combinatorics | Enumerative combinatorics is an area of combinatorics that deals with the number of ways that certain patterns can be formed. Two examples of this type of problem are counting combinations and counting permutations. More generally, given an infinite collection of finite sets Si indexed by the natural numbers, enumerative combinatorics seeks to describe a counting function which counts the number of objects in Sn for each n. Although counting the number of elements in a set is a rather broad mathematical problem, many of the problems that arise in applications have a relatively simple combinatorial description. The twelvefold way provides a unified framework for counting permutations, combinations and partitions.
The simplest such functions are closed formulas, which can be expressed as a composition of elementary functions such as factorials, powers, and so on. For instance, as shown below, the number of different possible orderings of a deck of n cards is f(n) = n!. The problem of finding a closed formula is known as algebraic enumeration, and frequently involves deriving a recurrence relation or generating function and using this to arrive at the desired closed form.
Often, a complicated closed formula yields little insight into the behavior of the counting function as the number of counted objects grows.
In these cases, a simple asymptotic approximation may be preferable. A function is an asymptotic approximation to if as . In this case, we write
Generating functions
Generating functions are used to describe families of combinatorial objects. Let denote the family of objects and let F(x) be its generating function. Then
where denotes the number of combinatorial objects of size n. The number of combinatorial objects of size n is therefore given by the coefficient of . Some common operation on families of combinatorial objects and its effect on the generating function will now be developed.
The exponential generating function is also sometimes used. In this case it would have the form
Once determined, the generating function yields the information given by the previous approaches. In addition, the various natural operations on generating functions such as addition, multiplication, differentiation, etc., have a combinatorial significance; this allows one to extend results from one combinatorial problem in order to solve others.
Union
Given two combinatorial families, and with generating functions F(x) and G(x) respectively, the disjoint union of the two families () has generating function F(x) + G(x).
Pairs
For two combinatorial families as above the Cartesian product (pair) of the two families () has generating function F(x)G(x).
Sequences
A (finite) sequence generalizes the idea of the pair as defined above. Sequences are arbitrary Cartesian products of a combinatorial object with itself. Formally:
To put the above in words: An empty sequence or a sequence of one element or a sequence of two elements or a sequence of three elements, etc.
The generating function would be:
Combinatorial structures
The above operations can now be used to enumerate common combinatorial objects including trees (binary and plane), Dyck paths and cycles. A combinatorial structure is composed of atoms. For example, with trees the atoms would be the nodes. The atoms which compose the object can either be labeled or unlabeled. Unlabeled atoms are indistinguishable from each other, while labelled atoms are distinct. Therefore, for a combinatorial object consisting of labeled atoms a new object can be formed by simply swapping two or more atoms.
Binary and plane trees
Binary and plane trees are examples of an unlabeled combinatorial structure. Trees consist of nodes linked by edges in such a way that there are no cycles. There is generally a node called the root, which has no parent node. In plane trees each node can have an arbitrary number of children. In binary trees, a special case of plane trees, each node can have either two or no children. Let denote the family of all plane trees. Then this family can be recursively defined as follows:
In this case represents the family of objects consisting of one node. This has generating function x. Let P(x) denote the generating function .
Putting the above description in words: A plane tree consists of a node to which is attached an arbitrary number of subtrees, each of which is also a plane tree. Using the operation on families of combinatorial structures developed earlier, this translates to a recursive generating function:
After solving for P(x):
An explicit formula for the number of plane trees of size n can now be determined by extracting the coefficient of xn:
Note: The notation [xn] f(x) refers to the coefficient of xn in f(x).
The series expansion of the square root is based on Newton's generalization of the binomial theorem. To get from the fourth to fifth line manipulations using the generalized binomial coefficient is needed.
The expression on the last line is equal to the (n − 1)st Catalan number. Therefore, pn = cn−1.
| Mathematics | Combinatorics | null |
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